Practical Astrodynamics (Gnv64)

Springer Aerospace Technology
Alessandro de Iaco Veris
Practical
Astrodynamics
More information about this series at http://www.springer.com/series/8613
Practical Astrodynamics
123
Rome
Italy
ISSN 1869-1730 ISSN 1869-1749 (electronic)

ISBN 978-3-319-62219-4 ISBN 978-3-319-62220-0 (eBook)
https://doi.org/10.1007/978-3-319-62220-0
Library of Congress Control Number: 2017945255
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This book is dedicated to the memory of
professor Luigi Broglio.
Preface
Ingegnati, se puoi, d'esser palese
(Dante, Vita nova, XIX).
The principle which has guided me in writing this book is indicated above. It has
been written for the purpose of making astrodynamics, and such mathematical
methods as are necessary to solve problems arising in this matter accessible to the
largest possible number of readers, and in particular to students, aerospace engi-
neers, and scientists.
To this end, the formal and sometimes esoteric language, which is still used in
some books dealing with the same matter, has been abandoned. On the contrary, all
possible means have been used to direct and keep the attention of the reader to the
subject. The presentation of theoretical concepts has been followed by their
application to practical cases. The numerical examples have been carried out from
beginning to end, in order to guide those who make their first attempts to apply the
methods described to cases of their interest. Equations and formulae appearing in
the text have been rewritten as many times as necessary, instead of identifying them
by numerals, which would have compelled the reader to search the pages in which
they had been presented the first time. No exercises have been proposed without a
complete guide to solve them. Several numerical tables have been given to enable
the reader to apply immediately the methods presented. Some items of the so-called
grey literature (unpublished works, technical reports, Ph.D. theses, lecture notes,
etc.) have also been consulted and cited. In some cases, the opinions expressed by
various authors who have previously written on the matter have been cited and
compared. Care has been taken not to lose sight of the practical problem which has
given rise to the mathematical method used for its solution. My goal has been rather
the search for a clear and simple manner of expressing the concepts to be introduced
to my readers than a strict observance of the classical form. Such are, if not the
results obtained, at least the means by which I have striven to attain my object. It is
to be hoped that these efforts will be favourably received by my readers.
vii
viii Preface
I have read carefully the manuscript of this book before printing. However,
I cannot believe it to be wholly free from any possible errors. Therefore, I shall be
grateful to any person who will call my attention on such errors as he or she may
discover.
Rome, Italy Alessandro de Iaco Veris

June 2016
Contents
1 The Two-Body Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1

1.1 Position of the Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1
1.2 The Conic Sections and Their Geometrical Properties . . . . . . . 12
1.3 The Elliptic Orbits . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 20
1.4 The Hyperbolic and Parabolic Trajectories . . . . . . . . . . . . . . . . 30
1.5 The Lambert Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 37
1.6 Transfer Times for Elliptic, Parabolic, and Hyperbolic
Trajectories . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 46
1.7 A Unified Form of Lambert’s Equations. . . . . . . . . . . . . . . . . . 58
1.8 An Example of Solution of Lambert’s Problem Using
Universal Variables . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 65
1.9 The Classical Orbital Elements . . . . . . . . . . . . . . . . . . . . . . . . . 75
1.10 Orbital Elements Defined for Any Orbit . . . . . . . . . . . . . . . . . . 89
1.11 The Lagrangian Coefficients f, g, f′, and g′ in Closed Form . . . 92
1.12 The Lagrangian Coefficients f and g in Time Series . . . . . . . . . 101
1.13 Canonical Units . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 105
1.14 The n-Body Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 108
1.15 The Halo Orbits . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 122
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 126
2 Orbit Determination from Observations . . . . . . . . . . . . . . . . . . . . . 129
2.2 Topocentric Co-ordinate Systems . . . . . . . . . . . . . . . . . . . . . . . 134
2.3 Orbit Determination from a Single Radar Observation . . . . . . . 137
2.4 The Measurement of Time in Astronomy . . . . . . . . . . . . . . . . . 149
2.5 Orbital Elements from Angle and Range Measurements . . . . . . 163
2.6 Orbital Elements from Three Measurements of Angles
(Method of Gauss) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 178
(Method of Laplace) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 212
ix
x Contents
2.8 Improvement in Orbit Determination by Differential

Correction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 222
2.9 Improvement in Orbit Determination by Weighted Least-
Squares Estimation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 225
2.10 Numerical Solution of the Least-Squares Estimation
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 239
2.11 The Kalman Filter . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 273
2.12 Numerical Methods for Kalman Filtering . . . . . . . . . . . . . . . . . 285
2.13 The Unscented Kalman Filter . . . . . . . . . . . . . . . . . . . . . . . . . . 299
2.14 The Square-Root Unscented Kalman Filter . . . . . . . . . . . . . . . . 303
2.15 The Minimax Filter. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 309
2.16 A More Robust Unscented Kalman Filter . . . . . . . . . . . . . . . . . 314
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 317
3 The Central Gravitational Force and Its Perturbations . . . . . . . . . 323
3.1 The System of Forces Acting on an Earth Satellite . . . . . . . . . 323
3.2 The Perturbation Due to the Non-spherical Earth . . . . . . . . . . . 325
3.3 The Changes of Orientation of the Earth Axis . . . . . . . . . . . . . 356
3.4 The Change of Co-ordinates Due to Precession . . . . . . . . . . . . 361
3.5 The Change of Co-ordinates Due to Nutation . . . . . . . . . . . . . . 365
3.6 The Change of Co-ordinates Due to the Rotation of the
Earth . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 374
3.7 The Change of Co-ordinates Due to Polar Motion . . . . . . . . . . 377
3.8 The Fundamental Reference Systems . . . . . . . . . . . . . . . . . . . . 379
3.9 The Frame-Bias Matrix . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 381
3.10 The Co-ordinate Transformation, Based on the Equinox,
Between the Celestial and Terrestrial Reference Systems . . . . . 384
3.11 The Co-ordinate Transformation, Based on the Non-rotating
Origins, Between the Celestial and Terrestrial Reference
Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 386
3.12 The Co-ordinate Transformation, According to the GOCE
Standards, Between the Celestial and Terrestrial Reference
Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 397
3.13 The Luni-Solar Perturbation . . . . . . . . . . . . . . . . . . . . . . . . . . . 407
3.14 The Position of the Perturbing Body . . . . . . . . . . . . . . . . . . . . 421
3.15 The Position of the Perturbing Body from NASA/JPL
Ephemeris Files . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 435
3.16 The Radiation Pressure Due to the Sun . . . . . . . . . . . . . . . . . . 438
3.17 The Eclipse Factor . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 454
3.18 The Radiation Pressure Due to the Earth . . . . . . . . . . . . . . . . . 478
3.19 The Atmospheric Drag . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 488
3.20 The Lifetime of an Earth Satellite Subject to
Atmospheric Drag . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 497
3.21 The Fundamental Properties of the Earth Atmosphere . . . . . . . 517
Contents xi
3.22 Atmospheric Density Models . . . . . . . . . . . . . . . . . . ........ 522

3.23 The Angular Velocity of the Atmosphere . . . . . . . . . ........ 541
3.24 The Relativistic Perturbations . . . . . . . . . . . . . . . . . . ........ 543
3.25 The Perturbations Due to Continuous Low-Thrust
Propulsion. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ........ 548
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ........ 559
4 Impulsive Orbital Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 571
4.2 Engines and Propellants for High-Thrust Rockets . . . . . . . . . . . 574
4.3 Launch Windows . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 611
4.4 Range Safety . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 619
4.5 Ascent Trajectories . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 624
4.6 Insertion into Orbit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 630
4.7 Rendezvous Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 635
4.8 Rendezvous-Compatible Orbits . . . . . . . . . . . . . . . . . . . . . . . . . 647
4.9 Intermediate Orbits for Rendezvous . . . . . . . . . . . . . . . . . . . . . 653
4.10 The Hill–Clohessy–Wiltshire Equations . . . . . . . . . . . . . . . . . . 668
4.11 The Hill–Clohessy–Wiltshire Equations Applied to
Rendezvous Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 676
4.12 Hohmann Transfer Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . 684
4.13 Bi-Elliptic Transfer Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . 697
4.14 Change of Orbital Plane . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 701
4.15 Change of the Position of a Spacecraft in Its Orbit . . . . . . . . . 708
4.16 Change of the Apsidal Line of an Orbit . . . . . . . . . . . . . . . . . . 711
4.17 Drag Make-up Manoeuvres for Satellites in Low-Altitude
Orbits . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 720
4.18 Manoeuvres for Geostationary Satellites . . . . . . . . . . . . . . . . . . 722
4.19 De-orbiting Manoeuvres . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 739
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 753
5 Interplanetary Trajectories . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 759
5.2 The Hohmann Ellipse Approximation . . . . . . . . . . . . . . . . . . . . 760
5.3 The Departure and Arrival Times . . . . . . . . . . . . . . . . . . . . . . . 763
5.4 The Spheres of Influence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 769
5.5 The Patched-Conic Approximation . . . . . . . . . . . . . . . . . . . . . . 773
5.6 The Departure of a Spacecraft from a Planet . . . . . . . . . . . . . . 774
5.7 The Arrival of a Spacecraft at a Planet . . . . . . . . . . . . . . . . . . . 787
5.8 The Flight of a Spacecraft Past a Planet . . . . . . . . . . . . . . . . . . 797
5.9 The Gravity Assist . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 811
5.10 Orbital Elements of the Planets. . . . . . . . . . . . . . . . . . . . . . . . . 815
5.11 General Interplanetary Trajectories . . . . . . . . . . . . . . . . . . . . . . 822
xii Contents
5.12 The Aerodynamic Assist . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 838

5.13 Trajectories of Vehicles Propelled by Solar Radiation
Pressure . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 841
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 854
6 Numerical Integration of the Equations of Motion . . . . . . . . . . . . . 857
6.2 Fundamental Concepts on the Runge–Kutta Methods . . . . . . . . 861
6.3 Runge–Kutta Fourth-Order Methods with Local Truncation
Error Control . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 866
6.4 Runge–Kutta Methods with Order Higher Than Four . . . . . . . . 871
6.5 Runge–Kutta–Nyström Methods . . . . . . . . . . . . . . . . . . . . . . . . 892
6.6 Step-Size Control with Runge–Kutta–Nyström Methods. . . . . . 898
6.7 Special Runge–Kutta Methods . . . . . . . . . . . . . . . . . . . . . . . . . 903
6.8 Special Runge–Kutta–Nyström Methods . . . . . . . . . . . . . . . . . . 910
6.9 Interpolants . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 924
6.10 Symplectic Explicit Special Nyström Methods . . . . . . . . . . . . . 934
6.11 Performance Comparison for Runge–Kutta(–Nyström)
Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 939
6.12 Bulirsch-Stoer Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 944
6.13 Multi-step Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 949
6.14 The Adams Method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 953
6.15 The Störmer-Cowell Method . . . . . . . . . . . . . . . . . . . . . . . . . . 968
6.16 The Gauss-Jackson Method . . . . . . . . . . . . . . . . . . . . . . . . . . . 974
6.17 Calculation of the Starting Values . . . . . . . . . . . . . . . . . . . . . . 982
6.18 Halving the Step Size . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 989
6.19 Integration for Elliptic Orbits of High Eccentricity . . . . . . . . . . 994
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 998
7 Dynamics of Rigid Bodies . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1001
7.1 The Motion of Rigid Bodies . . . . . . . . . . . . . . . . . . . . . . . . . . . 1001
7.2 The Matrix of Inertia . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1003
7.3 Kinetic Energy of a Rigid Body . . . . . . . . . . . . . . . . . . . . . . . . 1005
7.4 Moment of Inertia of a Rigid Body About an Arbitrary
Axis . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1006
7.5 Principal Axes of Inertia . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1008
7.6 Euler’s Equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1009
7.7 An Axially Symmetric (I1 = I2) Rotating Body not Subject to
External Moments . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1010
7.8 An Axially Symmetric (I1 = I2) Rotating Body not Subject to
External Moments (in Terms of Euler’s Angles) . . . . . . . . . . . 1013
7.9 Unsymmetrical Body Not Subject to External Moments
(Geometric Solution). . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1016
7.10 Unsymmetrical Body Not Subject to External Moments
(Analytic Solution) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1022
Contents xiii
7.11 Elementary Concepts on Elliptic Integrals . . . . . . . . . . . . . . . . 1025

7.12 Stability of the Rotation of a Rigid Body About Its Principal
Axes . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1031
7.13 General Motion of a Rigid Body . . . . . . . . . . . . . . . . . . . . . . . 1033
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1040
8 Instruments for Aerospace Navigation . . . . . . . . . . . . . . . . . . . . . . . 1041
8.1 Motion of a Symmetric Gyroscope . . . . . . . . . . . . . . . . . . . . . . 1041
8.2 Steady Precession of a Symmetric Gyroscope . . . . . . . . . . . . . 1049
8.3 Precession and Nutation of the Polar Axis of the Earth . . . . . . 1054
8.4 Small Oscillations of Gyroscopes . . . . . . . . . . . . . . . . . . . . . . . 1062
8.5 Oscillations of Gyroscopes About Gimbal Axes . . . . . . . . . . . . 1066
8.6 Effects Due to the Moments of Inertia of the Gimbals . . . . . . . 1072
8.7 The Gyrocompass . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1080
8.8 The Rate Gyroscope . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1092
8.9 The Rate Integrating Gyroscope . . . . . . . . . . . . . . . . . . . . . . . . 1096
8.10 High-Precision Gyroscopes . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1097
8.11 Optical Gyroscopes . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1099
8.12 Vibrating Structure Gyroscopes . . . . . . . . . . . . . . . . . . . . . . . . 1104
8.13 Accelerometers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1107
8.14 The Stable Platform . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1112
8.15 Inertial Navigation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1121
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1130
9 Attitude Stabilisation and Control of Earth Satellites . . . . . . . . . . . 1133
9.1 Attitude of Earth Satellites . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1133
9.2 Moments Due to Aerodynamic Forces . . . . . . . . . . . . . . . . . . . 1139
9.3 Moments Due to Electromagnetic Induction . . . . . . . . . . . . . . . 1141
9.4 Moments Due to Solar Radiation Pressure . . . . . . . . . . . . . . . . 1142
9.5 Moments Due to Gravity Gradient . . . . . . . . . . . . . . . . . . . . . . 1143
9.6 Moments Due to Micrometeorites . . . . . . . . . . . . . . . . . . . . . . . 1150
9.7 Comparison of the Magnitudes of the External Moments . . . . . 1150
9.8 Single-Spin and Dual-Spin Stabilisation of Satellites . . . . . . . . 1151
9.9 Nutation Dampers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1163
9.10 Gravity-Gradient Stabilisation of Satellites . . . . . . . . . . . . . . . . 1178
9.11 Stabilisation of Satellites by Means of Electromagnetic
Induction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 1183
9.12 Stabilisation of Satellites by Means of Reaction Jets . . . . . ... 1186
9.13 Stabilisation of Satellites by Means of Reaction Flywheels ... 1190
9.14 Stabilisation of Satellites by Means of Control Moment
Gyroscopes . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . ... 1194
xiv Contents
9.15 Three-Axis Controlled Satellites . . . . . . . . . . . . . . . . . . . . . . . . 1203

9.16 Attitude Re-orientation of a Satellite by Means of Impulse
Coning . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1205
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1208
10 Dynamics of Spinning Rockets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1213
10.1 The Motion of a Spinning Rocket . . . . . . . . . . . . . . . . . . . . . . 1213
10.2 Misalignment of the Thrust Vector in Body-Fixed Co-
ordinates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1214
10.3 Misalignment of the Thrust Vector in Inertial Co-ordinates . . . 1216
10.4 Near-Symmetric Body of Revolution Not Subject to
Moments . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1219
10.5 Rockets of Variable Mass . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1226
10.6 Damping Effect of the Exhaust Gas in a Non-spinning Rocket
of Variable Mass . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1228
10.7 Euler’s Equations for Spinning Rockets of Variable Mass . . . . 1229
10.8 Angle of Attack of a Rocket . . . . . . . . . . . . . . . . . . . . . . . . . . . 1234
10.9 The Motion of a Spinning Rocket with Varying Configuration
and Mass . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1238
10.10 The Yo-Yo de-Spin Mechanism . . . . . . . . . . . . . . . . . . . . . . . . 1244
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1255
11 Performance and Optimisation of Rockets . . . . . . . . . . . . . . . . . . . 1257
11.1 Performance of a Single-Stage Rocket . . . . . . . . . . . . . . . . . . . 1257
11.2 Multi-stage Rockets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1265
11.3 Optimum Staging for Multi-stage Rockets . . . . . . . . . . . . . . . . 1268
11.4 Optimum Trajectory to Place a Satellite into Orbit . . . . . . . . . . 1282
11.5 Optimum Consumption of Propellant . . . . . . . . . . . . . . . . . . . . 1289
11.6 Gravity Turn Trajectories . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1295
11.7 Trajectories of Long-Range Ballistic Missiles . . . . . . . . . . . . . . 1302
References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1308
Chapter 1
The Two-Body Problem
1.1 Position of the Problem
Astrodynamics is defined by Kaplan [1] as “the study of controlled flight paths of

man-made spacecraft”. This discipline is also defined in a more detailed manner by
the American Institute of Aeronautics and Astronautics [2] as “the determination,
prediction, physical adjustment, and optimisation of trajectories in space; space
navigation and mission analysis; perturbation theories and expansions; spacecraft
attitude dynamics and estimation”. These topics and the mathematical methods used
to solve practical problems arising in them form the subject of the present book.
Some authors separate astrodynamics from celestial mechanics, by limiting the
scope of the latter to the motion of natural celestial bodies, which is not under
human control. However, the motion of all celestial bodies, be they natural or
artificial, is governed by the same laws of mechanics.
This introductory chapter is meant to provide the reader with the basic concepts
of the two-body problem, which consists in determining the motion of two isolated
bodies, of masses, respectively, m1 and m2, attracting each other with Newtonian
gravitational forces whose magnitude is directly proportional to the product m1m2
and inversely proportional to the square of the distance between the two bodies. In
other words, given at an initial time the positions and velocities of two isolated
bodies acted upon only by their mutual gravitational attraction, it is required to
determine their positions and velocities at any other time.
As is well known, the solution of this problem is governed by Kepler’s laws of
planetary motion, which are given below.
© Springer International Publishing AG 2018 1

A. de Iaco Veris, Practical Astrodynamics, Springer Aerospace Technology,
https://doi.org/10.1007/978-3-319-62220-0_1
2 1 The Two-Body Problem
1. The orbits of the planets around the Sun are ellipses, with the Sun at one focus.
2. The radius vector of each planet (i.e., the straight line segment drawn from the
Sun to that planet) sweeps out equal areas in equal times, as the planet travels
along its orbit.
In the figure shown above, the time tAB taken by a planet to go from A to B is
the same as the time tCD taken by the same planet to go from C to D, because the
area of the sector FAB is equal to the area of the sector FCD. Thus, a planet
moves fastest along its orbit when it is in the region about perihelion, and most
slowly when it is near aphelion. For the sake of clearness, the figures given
above show an ellipse having an eccentricity of about 0.68, which value is much
higher than that of the orbit of any planet revolving around the Sun. The real
eccentricities of the planetary orbits are generally small, as shown in the fol-
lowing table.
1.1 Position of the Problem 3
Planetary eccentricities
Mercury 0.205 Jupiter 0.049
Venus 0.007 Saturn 0.057
Earth 0.017 Uranus 0.046
Mars 0.094 Neptune 0.011
Courtesy of NASA [3]
3. The ratio of the squares of the periods T1 and T2 of revolution for two planets is
equal to the ratio of the cubes of their respective major semi-axes a1 and a2, that is,
T12 a31
¼
T22 a32
As will be shown in Sect. 1.3, the third law implies that the period, T, of a
planet, which revolves around the Sun along an elliptic orbit, depends only on the
size of the major semi-axis, a, of the ellipse, that is,
3 12
a
T ¼ 2p
l
where l = GM = 1.327 1011 km3/s2 is the gravitational parameter of the Sun,

G = 6.673 10−20 km3/(kg s2) is the universal gravitational constant, and M is the
mass of the Sun. This also means that the mean orbital speed is lower for the planets
orbiting far away from the Sun than for those orbiting near it, as shown in the
following table.
Mean orbital speeds (km/s) of the planets
Mercury 47.4 Jupiter 13.1
Venus 35.0 Saturn 9.7
Earth 29.8 Uranus 6.8
Mars 24.1 Neptune 5.4
Courtesy of NASA [3]
Johannes Kepler found these laws (the first two in 1609 and the third in 1619) in
an effort to fit various geometrical curves to the astronomical observations per-
formed by Tycho Brahe on the orbit of Mars. In 1609, Kepler found that the data
supplied by Brahe could be explained if the path followed by Mars were an ellipse
having the Sun at one focus. The two-body model used by Kepler to describe the
motion of the planets around the Sun is valid, because the gravitational attraction
exerted by the Sun on the planets overcomes all other forces acting on them by
several orders of magnitude.
As has been shown at length by several authors (see for example Hyman [4]),
Kepler’s laws are a consequence of the general laws of motion discovered in 1687
by Isaac Newton. These laws are [5]:
1. Corpus omne perseverare in statu suo quiescendi vel movendi uniformiter in
directum, nisi quatenus a viribus impressis cogitur statum illum mutare (Every
body remains in its state of rest or of uniform motion in a straight line, unless it
is compelled to change that state by forces impressed upon it). This law was
originally enunciated by Galileo in 1638 [6] in the following terms “… Inoltre, è
lecito aspettarsi che, qualunque grado di velocità si trovi in un mobile, gli sia per
sua natura indelebilmente impresso, purché siano tolte le cause esterne di
accelerazione o di ritardamento; il che accade soltanto nel piano orizzontale …”
and is also known as the law of inertia.
2. Mutationem motus proportionalem esse vi motrici impressæ, & fieri secundum
lineam rectam qua vis illa imprimitur (The rate of change of momentum is
proportional to the motive force impressed and is directed along the straight line
in which that force is impressed), that is,
dð m v Þ
f ¼
dt
where m is the mass of the body considered, and f and v are, respectively, the
resultant of the forces applied to the body and its velocity vector. In case of
constant mass, the preceding expression reduces to
f ¼ ma
where a = dv/dt is the acceleration vector applied to the body.

3. Actioni contrariam semper & æqualem esse reactionem: sive corporum duorum
actiones in se mutuo semper esse æquales & in partes contrarias dirigi (To every
action there is always opposed an equal reaction, that is, the mutual actions
exerted by two bodies are always equal and oppositely directed).
To show how Kepler’s laws on planetary orbits follow from Newton’s laws, let
us make the following assumptions for the two-body problem:
• each of the two bodies can be considered as a particle having all its mass
concentrated in its centre of mass; and
• each of the two bodies is only subject to the gravitational force produced by the
other body and directed along the straight line joining the centres of mass of the
two bodies.
With these assumptions, we study the motion of a body in an inertial (i.e.,
non-accelerated and non-rotating) reference system. Following Kaplan [1], let m1
and m2 be the masses of the two isolated bodies. Let XYZ be a Cartesian inertial
system having its origin in O. Let C be the centre of mass of the two bodies. Let
r1 ¼ X1 uX þ Y1 uY þ Z1 uZ
r2 ¼ X2 uX þ Y2 uY þ Z2 uZ
rC ¼ XC uX þ YC uY þ ZC uZ
be the position vectors of, respectively, m1, m2, and C, where uX, uY, and uZ are the
unit vectors along, respectively, X, Y, and Z. Let r be a vector defined as follows
r ¼ r1 r2
and shown in the following figure.
Since C is, by definition, the centre of mass of m1 and m2, then the position
vector of C is
m1 r1 þ m2 r2
rC ¼
m1 þ m2
that is,

m2
r1 rC ¼ r
m1 þ m2
This is because, by definition, r = r1 − r2. Likewise, there results

m1
r2 rC ¼ r
m1 þ m2
Let f1 and f2 be the forces acting on, respectively, m1 and m2. Then, Newton’s
second law of motion, written for m1 and m2, states that

m1 m2
f 1 ¼ m 1 a1 ¼ m1 r001
¼ r00 þ m1 r00C
m1 þ m2

00 m1 m2
f 2 ¼ m2 a2 ¼ m2 r2 ¼ r00 þ m2 r00C
m1 þ m2
where a double prime sign (″) indicates second derivatives with respect to time.
Since f1 and f2 are forces of mutual attraction, then f1 + f2 = 0. Thus, there
results
m1 r00C þ m2 r00C ¼ 0
that is,
r00C ¼ 0
This means that the centre of mass of the two bodies does not accelerate.
Substituting r00C ¼ 0 into the equations written above

m1 m2
f 1 ¼ m1 a1 ¼ m1 r001 ¼ r00 þ m1 r00C
m1 þ m2

m1 m2
f 2 ¼ m2 a2 ¼ m2 r002 ¼ r00 þ m2 r00C
m1 þ m2
yields

m1 m2
f 1 ¼ m1 a1 ¼ m1 r001 ¼ r00
m1 þ m2

m1 m2
f 2 ¼ m2 a2 ¼ m2 r002 ¼ r00
m1 þ m2
The attracting force f1 is a Newtonian gravitational force, such that

m1 m2 r
f 1 ¼ G
r2 r
where the minus sign is due to the attracting nature of f1, which acts along r but
oppositely to r. After setting l = G(m1 + m2), the differential equation governing
the motion for the two-body problem can be written as follows
l
r00 þ r¼0
r3
This equation can be solved analytically, because some vector operations exe-
cuted on it give rise to exact differentials, which are immediately integrable, as will
be shown below. The constants of integration are also called orbital elements. The
vector product of r by the two members of this equation yields
l
r r00 þ r r¼0
r3
Because of the identity r r 0, the preceding expression reduces to
r r00 ¼ 0
Since (r r′)′ = r′ r′ + r r″ = r r″, and since r′ r′ = 0, r r″ = 0,

where a prime sign (′) denotes first derivatives with respect to time, then
ðr r0 Þ0 ¼ 0
Thus, h = r r′, which is the moment of momentum per unit mass, does not
vary with time during the motion. The result h = constant provides three integrals
of motion. Since h = r r′ is orthogonal to both the position (r) and velocity (r′)
vectors which define the plane of motion, then the constancy of h implies the
constancy of the plane of motion. In other words, the motion is confined to an
invariable plane, which is orthogonal to the constant vector h.
In addition, since h is constant not only in direction, but also in magnitude, then
the areal velocity of a planet revolving around the Sun is also constant, as shown in
the preceding figure.
Let P be the position, at a time t, of a given planet revolving about the Sun. Let
r and / be the polar co-ordinates of P with respect to a reference system having its
origin in the centre of mass, S, of the Sun. Let Q be the position of the planet along
its orbit after the lapse of an infinitesimal interval of time dt.
During this interval, the infinitesimal area swept out by the radius vector r is
1 1
AreaðPSQÞ ¼ r ðr d/Þ ¼ r 2 d/
2 2
Thus, the instantaneous areal velocity of the planet is
d½AreaðPSQÞ 1 2 d/ 1 2 0 1 2
¼ r r / ¼ r x
dt 2 dt 2 2
where x is the instantaneous angular velocity of the planet. Since the magnitude
(h = r2x = r2/′) of the moment of momentum per unit mass of the planet is
constant, as has been shown above, then the areal velocity of the planet will also be
constant during its time of revolution. This proves Kepler’s second law.
Now, the vector product of the two members of the equation
l
r00 þ r¼0
r3
by h yields
l l
r00 h ¼ r h ¼ r ðr r0 Þ
r3 r3
By applying the following identity (rule of triple vector products)
a ðb cÞ ¼ bða cÞ cða bÞ
to the expression −(l/r3)r (r r′), there results

l l l
r ðr r0 Þ ¼ ½ r ðr r0 Þ ¼ ½rðr r0 Þ r0 ðr rÞ
r3 r3 r3
Since r r′ = rr′ and r r = r2, then

l l
3 ½rðr r0 Þ r0 ðr rÞ ¼ 3 r rr 0 r0 r 2 ¼ lðr=r Þ0
r r
that is,
r00 h ¼ lðr=r Þ0
Since h is a constant vector, then the preceding differential equation can be

integrated immediately, with the following result
r
r0 h ¼ l þe
r
where e is a constant of integration called the eccentricity vector. The preceding

expression provides three further integrals of motion.
The vector e lies in the orbit plane, because the dot product of (r′ h) by h,
that is,
ðr0 hÞ h
is equal to 0. This is because (r′ h) is orthogonal to h. Therefore, the direction of

e can be taken as a reference direction in the orbital plane.
The dot product of the expression r′ h = l(r/r + e) by r yields
h r i
ðr0 hÞ r ¼ l þ e r
r
As to the left-hand side of the preceding equality, the two following identities
a ðb cÞ ¼ b ðc aÞ ¼ c ða bÞ
ab¼ba
yield
ðr0 hÞ r ¼ r ðr0 hÞ ¼ h ðr r0 Þ ¼ h h ¼ h2
As to the right-hand side of the same equality, there results

h r i l
l þe r ¼ ðr rÞ þ l er cos / ¼ lr ð1 þ e cos /Þ
r r
where /, called true anomaly, is the angle between the position vector, r, and the
eccentricity vector, e, that is,
er
cos / ¼
er
Thus, the equality can also be written as follows
h2 ¼ lrð1 þ e cos /Þ
which, solved for r, yields
h2
l
r¼
1 þ e cos /
The expression written above is the equation of a conic section in polar

co-ordinates r and /, where the origin of the radius vector r is in a focus, and the
origin of the true anomaly / is the direction of the point of the conic section which
is nearest to that focus. This proves Kepler’s first law.
As has been shown above, the equation of a conic section in polar co-ordinates is
such that the value / = 0 of the true anomaly / corresponds to the minimum value
rp of the radius vector r. Consequently, the eccentricity vector e, which is always
contained in the plane of motion, is parallel to the direction of the minimum value
of r, that is, to the direction of perihelion.
By setting p = h2/l, where p has the dimension of length and is called semi-latus
rectum, the equation of a conic section in polar co-ordinates can also be written as
follows
p
r¼
1 þ e cos /
In this connexion, it is to be remembered that the semi-latus rectum of a conic

section is the distance from the single focus (if the conic section be a parabola), or
one of the two foci (if the conic section be an ellipse or a hyperbola), to the conic
section itself, measured along a line perpendicular to its major axis, as shown by the
preceding figure for a planetary orbit.
The differential equation derived above
l
r00 þ r¼0
r3
has important consequences for the mechanical energy possessed by a planet

moving around the Sun.
Let us consider the magnitude e of the eccentricity vector e. The square e2 of this
magnitude results from the scalar product of the vector e by itself, that is,
e2 ¼ e e
As has been shown above, there results

r
r0 h ¼ l þe
r
This equation, solved for e, yields
r0 h r
e¼
l r
Thus,
0 0
r h r r h r
ee¼
l r l r
1 0 2 h r i
¼ 2 ½ðr hÞ ðr0 hÞ ðr0 hÞ þ1
l l r
1 2
¼ 2 ½ðr0 hÞ ðr0 hÞ ðr r0 hÞ þ 1
l lr
Since (r′ h) (r′ h) = r′ h (r′ h) and since r′ and h are orthogonal

vectors, then
ðr0 hÞ ðr0 hÞ ¼ ðr 0 Þ2 h2 v2 h2
On the other hand,
r r0 h ¼ r r0 h ¼ h h ¼ h2
Thus,
v2 h2 2h2
e2 ¼ þ1
l2 lr
and
2 2 2
v h 2h2 2h2 v2 h2 h 2 v2
1 e2 ¼ 1 þ 1 ¼ ¼
l2 lr lr l2 l r l
The preceding equation, remembering that p = h2/l, becomes

2 v2
1 e2 ¼ p
r l
The factor (2/r − v2/l) must be a constant of the motion and must also be such
as to make the product p(2/r − v2/l) a dimensionless quantity. By defining
1
a¼
2 v2

r l
where the quantity a has the dimension of length, there results

p
1 e2 ¼
a
The previous expression, inserted into 1 − e2 = p(2/r − v2/l), leads to
1 2 l l
v ¼ constant ¼
2 r 2a
Since ½v2 and −l/r are, respectively, the kinetic energy per unit mass and the
potential energy per unit mass of a given planet revolving about the Sun, then the
constant E = −l/(2a) is the mechanical energy per unit mass of that planet. By the
way, the potential energy per unit mass of a planet revolving about the Sun is
negative, because gravity is an attracting force and the potential energy per unit
mass is defined to be zero for a planet at infinite distance from the Sun.
The preceding expression may also be written as follows

2 1
v ¼l
2
r a
which is the so-called energy integral or vis-viva integral.

The quantities p, a, and e are related by the equation shown above
p ¼ að1 e2 Þ
Since p is never negative, then e is less or greater than unity, depending on

whether a is positive or negative. The eccentricity, e, is equal to unity in case of
rectilinear motion (h = 0) or in case of zero mechanical energy (½v2 = l/r).
1.2 The Conic Sections and Their Geometrical Properties
As has been shown in Sect. 1.1, the solution of the differential equation governing
the two-body problem
l
r00 þ r¼0
r3
is a conic section having the following equation in polar co-ordinates (r, /)

1.2 The Conic Sections and Their Geometrical Properties 13
p
r¼
1 þ e cos /
The present section describes the three fundamental types of conic sections
(ellipse, parabola, and hyperbola) depending on the value of the eccentricity, e.
They are shown in the following figure, due to the courtesy of the Government of
the United States [7].
According to the most common definition, a conic section is a curve resulting

from the intersection of a plane with the surface of a circular right cone, as shown
on the left-hand side of the following figure.
The generation of a conic surface may be described as follows. Let us consider a

vertical straight line (called axis) and another straight line (called generatrix)
intersecting the axis at an angle a (called the vertex angle), as shown in the central
part of the preceding figure. The generatrix, in its rotation about the vertex, sweeps
out a surface, shown on the right-hand side, which is called a circular right cone, or
briefly a cone. Needless to say, both the upper and lower halves of the cone extend
indefinitely, respectively, upwards and downwards.
Now, let a cone, constructed as has been shown above, be intersected with a
plane. Depending on which value is chosen for the angle of intersection between
cone and plane, we may have the following cases:
• the plane intersecting the cone is orthogonal to the (vertical) axis, in which case
only one of the two halves of the cone is intersected, and the resulting closed
curve is a circumference (which degenerates to a point, in case of the inter-
section taking place just at the vertex); or
• the angle formed by the plane and the (vertical) axis is greater than the angle a
of the generatrix, in which case, again, the plane intersects only one of the two
halves of the cone, and the resulting closed curve is an ellipse; or
• the angle formed by the plane and the (vertical) axis is equal to the angle a of the
generatrix, in which case the plane is parallel to the side of the cone, and the
resulting open curve is a parabola; or
• the angle formed by the plane and the (vertical) axis is less than the angle a of
the generatrix, in which case the plane intersects both the halves of the cone, and
the resulting open curve, hyperbola, has two separate branches located in the
two halves of the cone.
Menaechmus (fourth century BC) seems to have been the first among the Greek
mathematicians to describe the conic sections [8, p. 46, and 9, p. 251].
According to a further definition, also due to the ancient Greeks, a conic section
is either a circumference or a locus of points P such that their absolute distance FP
from a fixed point F, called focus, is proportional to their absolute distance PP′ from
a fixed straight line, called directrix, as shown in the following figure.
Let e, called eccentricity, be the positive constant of proportionality. If e = 1,

then the conic section is a parabola. Otherwise (e 6¼ 1), let the major semi-axis, a,
be defined as follows
p
a¼
1 e2
where p is the semi-latus rectum, which is shown in the preceding figure.

If 0 < e < 1 and 0 < a < ∞, then the conic section is an ellipse; otherwise, if
e > 1 and −∞ < a < 0, then the conic section is a hyperbola.
The preceding figure illustrates this definition in case of an ellipse, whose
eccentricity e is less than unity (e 0.68 in the figure): an ellipse is the locus of
points P such that the ratio FP/PP′ has a constant value, which is the eccentricity.
The specific properties of the conic trajectories having, respectively, circular,
elliptic, parabolic, and hyperbolic shape are discussed below. The focus and the
eccentricity of a conic section have physical meaning as far as orbital motion is
concerned. The conic sections have generally two foci, F and F′, which coincide in
case of a circumference (e = 0). One of the foci, let it be F, is the location of the
central attracting body, whereas the other or vacant focus (F′) has no physical
meaning in orbital mechanics. In case of a parabola (e = 1 and a = ∞), the focus F′
is assumed to be located at an infinite distance to the left of the occupied focus F, as
shown in the following figure.
In case of a parabola (e = 1), the general equation of conic sections in polar

co-ordinates
p
r¼
1 þ e cos /
becomes
p
r¼
1 þ cos /
Since the following expression

/
1 tan2
2
cos / ¼
/
1 þ tan2
2
is a trigonometric identity, then
2
1 þ cos / ¼
/
1 þ tan2
2
and the equation of a parabola in polar co-ordinates can also be written as follows

1 /
r ¼ p 1 þ tan2
2 2
Since Sect. 1.1 has shown that
h ¼ r 2 /0 ¼ ðl pÞ2
1
then, substituting

1 2 /
r ¼ p 1 þ tan
2 2
into the preceding equation yields

2
1 2 2 / d/ 1
p 1 þ tan ¼ ð l pÞ 2
4 2 dt
By simplifying and separating the variables, the preceding differential equation

becomes
12 2
l 2 /
4 dt ¼ 1 þ tan d/
p3 2
In addition, since

/ 1
1 þ tan 2
¼ ¼ sec2 /
2 cos2 /
is a trigonometric identity, then the preceding differential equation becomes

12
l 4 /
4 3 dt ¼ sec d/
p 2
which, integrated, yields
12 Z/
l h
4 3 t tp ¼ sec4 dh
p 2
0
where tp is the time of perihelion passage. By defining a new variable w = h/2, the
preceding expression becomes
/
12 Z2
l
4 3 t tp ¼ 2 sec4 w dw
p
0
R
Thus, the problem reduces to computing sec4 w dw from 0 to //2. Since
Z Z Z Z
1 1 1 1
sec w dw ¼
4
dw ¼ dw ¼ dðtan wÞ
cos4 w cos2 w cos2 w cos2 w
then, integrating by parts, there results

Z Z
1 tan w 2 cos w sin w
dðtan wÞ ¼ ðtan wÞ dw
cos2 w cos2 w cos4 w
Z
tan w 1 cos2 w
¼ 2 dw
cos2 w cos4 w
Z Z
tan w dw dw
¼ 2 þ2
Z
tan w dw
¼ 2 þ 2 tan w
cos2 w cos4 w
Thus,
Z Z
dw tan w dw
¼ 2 þ 2 tan w
Z
dw tan w
3 ¼ þ 2 tan w
cos4 w cos2 w
Z
dw 1 tan w 1 2
4
¼ 2
þ 2 tan w ¼ sec2 w tan w þ tan w þ c
cos w 3 cos w 3 3
Now, since
sec2 w ¼ 1 þ tan2 w

Z
1
2 1
sec4 w dw ¼ 1 þ tan2 w tan w þ tan w þ c ¼ tan3 w þ tan w þ c
3 3 3
It follows that
/
12 Z2
l 2 / /
4 3 t tp ¼ 2 sec4 w dw ¼ tan3 þ 2 tan
p 3 2 2
0
By defining the quantity B = 3(l/p3)½ (t − tp), the preceding equation reduces to

/ /
tan3 þ 3 tan ¼ 2B
2 2
which is called Barker’s equation.

The following figure shows a hyperbola.
The reason, which accounts for the minus sign placed in front of the quantities
a (transverse semi-axis) and ea (semi-distance of the two foci) in the case of a
hyperbola, is the definition given above
p
a¼
1 e2
If e > 1, then a and ea must have negative values. As shown in the figure, on
each of the two sides of a hyperbola, the two arms of the curve approach indefi-
nitely two intersecting straight lines, called asymptotes of the hyperbola.
The quantity b for a hyperbola is the distance between FF′ (the straight line
joining the two foci) and the asymptote, measured orthogonally to FF′, that is,
1
b ¼ a e2 1 2
The angle, d, between the asymptotes is the bending of the curve from −∞ to
+∞ and has a physical meaning, because it represents the angle through which the
path of a spacecraft approaching a given planet is deviated by the force of gravi-
tational attraction exerted by the planet on it. This angle is called the flyby
turn-angle and depends on the eccentricity, e, of the hyperbola according to the
following definition

d a 1
sin ¼ ¼
2 ea e
1 1
d b aðe2 1Þ2 ðe2 1Þ2
cos ¼ ¼ ¼
2 ea ea e
which means that the greater the eccentricity of the hyperbola representing the path
of a spacecraft approaching a planet, the lesser the spacecraft deviates from its path
as a result of the force of gravity exerted by the planet.
It is to be borne in mind that only one of the two branches (i.e., either the
left-hand or the right-hand branch) of a hyperbola is a possible path for a spacecraft
approaching a planet, depending on which focus (either F or F′) is occupied by the
centre of mass of the planet approached. For example, if F is the occupied focus and
F′ is the vacant focus, then the left branch of the curve is the only possible path for
an approaching spacecraft. This is because of the attracting nature of the gravita-
tional force.
1.3 The Elliptic Orbits
As has been shown in Sect. 1.1, the orbits of the planets around the Sun and also
those of the artificial satellites around the Earth are ellipses (first Kepler’s law).
They are closed curves characterised by their period, which is the time necessary
for the orbiting body to make an entire revolution around the main attracting body.
In the following discussion, it will be assumed that such orbits are those of the
planets of the Solar System; however, the same considerations hold for any bodies
whose masses are much smaller than the mass of the principal attracting body.
1.3 The Elliptic Orbits 21
The elliptic orbit of a given planet about the Sun has a perihelion (or point of
smallest distance from the Sun) and an aphelion (or point of greatest distance from
the Sun). The same holds for an artificial satellite of the Earth, which has a perigee
(or point of smallest distance from the Earth) and an apogee (or point of greatest
distance from the Earth). Let rp and ra be, respectively, the radius vector of the
perihelion (or perigee) and the radius vector of the aphelion (or apogee), as shown
in the preceding figure. A simple inspection of this figure shows that
ra þ rp ¼ 2a
ra rp ¼ 2ea
where 2a is the length of the major axis of the elliptic orbit. Dividing the terms of
the second of these equalities by the terms of the first yields
ra rp
e¼
ra þ rp
To show how to determine the position P of a planet at a given time along its
orbit about the Sun (or the position of an artificial satellite along its orbit about the
Earth), we introduce the auxiliary variable, Æ, called eccentric anomaly.
This variable, shown in the following figure, is defined as follows
x ¼ aðcos Æ eÞ

y ¼ a 1 e2 sin Æ
Let us consider the elliptic orbit of a planet about the Sun, which is in the
occupied focus F. Let x, y be a Cartesian reference system having its origin in F.
Let P be the position of the planet at a given time t. Let r and / be, respectively,
the radius vector and the true anomaly of P at the same time t. The Cartesian
co-ordinates x and y of P are expressible as functions of the polar co-ordinates r and
/ as follows
x ¼ r cos /
y ¼ r sin /
Now, we draw an auxiliary circumference, whose centre, O, bisects F′F and

whose radius is equal to the major semi-axis, a, of the planetary orbit.
Starting from P, we produce a straight line PQ orthogonally to the line of apsides

(joining the two points of perihelion and aphelion), Q being the point of intersection
between PQ and the auxiliary circumference. Let Æ be the angle between OQ and
the line of apsides.
To show the relation existing between an elliptic orbit and its auxiliary cir-
cumference, we consider a Cartesian reference system x*, y* having its origin in O.
As is well known, the equations of the elliptic orbit and its auxiliary circumference,
with respect to Ox*y*, are, respectively,
x2 y2
þ 2 ¼1 ðellipse)
a2 b
x2 y2
þ 2 ¼1 ðcircumference)
a2 a
The preceding equations, solved for y*, yield
b 2 1
ye ¼ a x2 2
a
1
yc ¼ a2 x2 2
where the subscripts e and c identify the ordinates of, respectively, the elliptic orbit
and the associated auxiliary circumference.
Thus, the relation between the two sets of ordinates is
ye b
¼
yc a
Let us consider the following figure.
Following Bate et al. [10], we note that the area swept out by the radius vector r,
when the true anomaly varies from zero (at perihelion) to / (at the point P), is equal
to the difference
AreaðPSVÞ AreaðPSFÞ
Since the base SF of the triangle PSF is
SF ¼ OF OS ¼ ae a cos Æ
and the altitude SP of the same triangle is

b b
SP ¼ ye ¼ yc ¼ a sin Æ
a a
then there results

1 b
AreaðPSFÞ ¼ ðae a cos ÆÞ a sin Æ
2 a
1
¼ abðe sin Æ sin Æ cos ÆÞ
2
Let us consider now the area of the elliptic segment PSV. Since, as shown above,
the following equality holds
ye b
¼
yc a
then there will also result
b
AreaðPSVÞ ¼ AreaðQSVÞ
a
where QSV is the segment of the auxiliary circle corresponding to PSV.

Now, the area of the circular segment QSV is equal to
AreaðQSVÞ ¼ AreaðOVQÞ AreaðOSQÞ
These two areas, in turn, are
1
AreaðOVQÞ ¼ a2 Æ
2
(where Æ is measured in radians) and
1
AreaðOSQÞ ¼ ða cos Æ Þða sin Æ Þ
2
Thus,
1 1 1
AreaðQSVÞ ¼ a2 Æ ða cos Æ Þða sin Æ Þ ¼ a2 ðÆ sin Æ cos Æ Þ
2 2 2
hence
b 1
AreaðPSVÞ ¼ AreaðQSVÞ ¼ abðÆ sin Æ cos Æ Þ
a 2
Now, remembering that the area swept out by the radius vector r, when the true
anomaly varies from zero to /, is equal to the difference
AreaðPSVÞ AreaðPSFÞ
then the area swept out by the radius vector r, within the same interval of true
anomalies, is equal to
1 1
abðÆ sin Æ cos Æ Þ abðe sin Æ sin Æ cos Æ Þ
2 2
1
¼ abðÆ e sin Æ Þ
2
Since the area of the complete elliptic orbit is

b
pa2 ¼ pab
a
(where pa2 is the area of the auxiliary circle), then the interval of time taken by a
planet to go from its perihelion (whose true anomaly is zero) to a point P (whose
true anomaly is /) is

T 1 T
t tp ¼ abðÆ e sin Æ Þ ¼ ðÆ e sin Æ Þ
pab 2 2p
where tp is the time of perihelion passage for the given planet.

The period T of an elliptic orbit is
3 12
a
T ¼ 2p
l
where a is the major semi-axis of the planet considered, and l = GM is the grav-
itational parameter of the Sun. To show this, we remember that the areal velocity of
a planet along its orbit is constant, in accordance with Kepler’s second law. This
means that the area (pab) swept out by the radius vector, r, during a complete orbit
divided by the orbital period (T) is equal to ½h, where h is the magnitude of the
angular momentum per unit mass, that is,
pab 1 2 d / 1
¼ r ¼ h
T 2 dt 2
As has been shown in Sect. 1.1, there results
h2
p¼
l
where p is the semi-latus rectum of the planetary orbit, and

p
b¼ 1
ð1 e 2 Þ2
where b is the minor semi-axis of the orbit.

After squaring the equation pab/T = ½h and then introducing p = h2/l and
b = p/(1 − e2)½ into it, there results
p2 a2 b2 1 2
¼ h
2 T2 4
p 1 1
2 2
p a ¼ lp
1 e2 T 2 4
Dividing the two members of this equality by p yields

p 1 1
p2 a2 ¼ l
1e T
2 2 4
Since a = p/(1 − e2), then there results
p2 a3 1
¼ l
T2 4
which, solved for T, yields

3 12
a
T ¼ 2p
l
This proves Kepler’s third law. Substituting this expression into the equation
T
t tp ¼ ðÆ e sin Æ Þ
2p
yields
3 12
a
l
By setting
M ¼ Æ e sin Æ
where M is called the mean anomaly, and

l 12
n¼
a3
where n is called the mean motion, there results

M ¼ n t tp ¼ Æ e sin Æ
which is known as Kepler’s equation. A simple inspection of the preceding figure

shows that the equation of an ellipse given above, that is,
x2 y2
þ 2 ¼1
a2 b
can be expressed in terms of a, b, and Æ, as follows
x ¼ a cos Æ
y ¼ b sin Æ
The same figure also shows that
r ¼ a ae cosÆ
The preceding expression, compared with the polar equation of a conic section
að 1 e 2 Þ
r¼
1 þ e cos /
yields
e þ cos /
cosÆ ¼
1 þ e cos /
cosÆ e
cos / ¼
1 e cosÆ
In addition, since
1 að 1 e 2 Þ
y ¼ b sin Æ ¼ a 1 e2 2 sin Æ ¼ r sin/ ¼ sin/
1 þ e cos /
then there results

1
ð1 e2 Þ2 sin /
sin Æ ¼
1 þ e cos /
1
ð1 e2 Þ2 sin Æ
sin / ¼
1 e cos Æ
Finally, since

/ 1 cos/ / 1 þ cos/
sin2 ¼ cos2 ¼
2 2 2 2

Æ Æ Æ Æ
cos Æ ¼ cos2 sin2 1 ¼ cos2 þ sin2
2 2 2 2
are trigonometric identities, then there results

/ 1 cos Æ e að 1 þ e Þ 2 Æ
sin2 ¼ 1 ¼ sin
2 2 1 e cos Æ r 2

/ 1 cos Æ e að 1 e Þ 2 Æ
cos2 ¼ 1þ ¼ cos
2 2 1 e cos Æ r 2
Dividing the second expression by the first and taking the square root of the
quotient, there results
1
/ 1þe 2 Æ
tan ¼ tan
2 1e 2
What is known in practice is the true anomaly, /, of a planet with respect to its
position at perihelion (or of an artificial satellite with respect to its position at
perigee). Hence, we need a relation between / and the corresponding eccentric
anomaly, Æ. To this end, by inspection of the following figure
we note that OS = OF + FS, that is,
a cos Æ ¼ ae þ r cos /
Substituting
p að1 e2 Þ
r¼ ¼
1 þ e cos / 1 þ e cos /
into the preceding expression yields


að 1 e 2 Þ
a cos Æ ¼ ae þ cos /
1 þ e cos /
which in turn, simplified and solved for cos Æ, is
e þ cos /
cos Æ ¼
1 þ e cos /
The preceding expression makes it possible to compute the eccentric anomaly Æ

corresponding to a given true anomaly /. The quadrant where Æ actually lies is
determined, because / and Æ are always in the same half-plane with respect to the
line of apsides; that is, if 0 < / < p, then 0 < Æ < p; likewise, if p < / < 2p, then
p < Æ < 2p.
Bate et al. [10] illustrate this concept by means of the following example. With
reference to the following figure, it is required to determine the time taken by a
given planet to move along its orbit from a point A (identified, with respect to the
perihelion direction, by its true anomaly /A) to another point B (identified, with
respect to the same direction, by its true anomaly /B), if the planet does not pass
through its perihelion to move from A to B.
In this case (characterised by the fact that the true anomaly at the starting point is
smaller than that at the arrival point), the time taken by the planet to move from
A to B is

tB tA ¼ tB tp tA tp
where tp is the time of perihelion passage for that planet.

By contrast, it is required now to determine the time taken by the same planet to
move along its orbit from a point R (identified, with respect to the perihelion
direction, by its true anomaly /R) to another point S (identified, with respect to the
same direction, by its true anomaly /S), if the planet does pass through its peri-
helion to move from R to S, as shown in the following figure.
In the second case (characterised by the fact that the true anomaly at the starting
point is greater than that at the arrival point), the time required is

tR tS ¼ T þ tS tp tR tp
where tp, again, is the time of perihelion passage and T is the orbital period for the
planet considered.
In the general case, let ti and tf denote, respectively, the initial time (i.e., the time
at which a given planet starts from a given point) and the final time (i.e., time at
which the same planet arrives at another given point) along the planetary orbit. Let
Æi and Æf denote the correspondent eccentric anomalies.
The time interval Dt (tf − ti) necessary for the planet to move from the initial
point to the final point is
3 12
a

Dt ¼ 2mp þ Æ f e sin Æ f ðÆ i e sin Æ i Þ
l
where the integer m specifies how many times the planet passes through its peri-
helion to move from the initial point to the final point along its orbit.
1.4 The Hyperbolic and Parabolic Trajectories
The expressions derived in the previous section for the case of elliptic orbits have
correspondent expressions for the cases of hyperbolic and parabolic trajectories.
Following Battin [11], a hyperbolic trajectory can be expressed in parametric form
as follows
1.4 The Hyperbolic and Parabolic Trajectories 31
x ¼ a sec f
y ¼ b tan f
where the angle f, called the Gudermannian, is shown in the following figure.
Let us produce, from a point P of a given hyperbolic trajectory, a straight line PQ

orthogonal to the line of apsides (x*-axis) of the hyperbola, Q being the point where
this straight line intersects the line of apsides. From Q let us produce the tangent QR
to an auxiliary circumference having its centre O in the centre of the hyperbola and
its radius equal to −a (the transverse semi-axis of the hyperbola).
Let R be the point of tangency between the straight line produced from Q and the
auxiliary circumference. Let f be the angle which OR forms with the line of
apsides. This definition of f implies that
x2 y2 1 sin2 f 1 sin2 f

¼ ¼ ¼1
a2 b2 cos2 f cos2 f cos2 f
which is the well-known equation of a hyperbola in the Cartesian co-ordinates x*

and y* with origin in O.
The preceding figure shows that the distance OF = −ea is equal to the sum of
three terms, which are the projections of OR, RQ, and FP on the line of apsides:
OF ¼ OR cos f þ RQ sin f þ FP cos /

Now, OR is equal to −a; RQ is equal to −a tan f, because RQ is orthogonal to

OR; and FP is the radius vector r. Thus,
ea ¼ a cos f a tan f sin f þ r cos /
which simplified yields
ea ¼ a sec f þ r cos /
Consequently, the angle f used for hyperbolic trajectories has a geometric

analogy with the eccentric anomaly Æ used for elliptic orbits. In both cases, we use
auxiliary circumferences, each of them having its centre in the centre of the conic
section, be it an ellipse or a hyperbola, and its radius equal to, respectively, the
major semi-axis (a) of the elliptic orbit or the transverse semi-axis (−a) of the
hyperbolic trajectory. The formulae relating the angle f and the true anomaly / are
given below [12]:
1
e sec f ðe2 1Þ2 tan f
cos / ¼ sin/ ¼
e sec f 1 e sec f 1
1
e þ cos / ðe2 1Þ2 sin/
sec f ¼ tan f ¼
In addition, since

2 / að e þ 1Þ 2 f / aðe 1Þ 2 f
sin ¼ sin cos 2
¼ cos
2 r cos f 2 2 r cos f 2
then, by dividing the second expression by the first and taking the square root of the
quotient, there results
1
/ eþ1 2 f
tan ¼ tan
2 e1 2
The time equation for hyperbolic trajectories (corresponding to Kepler’s equa-

tion for elliptic orbits) is

f p
N ¼ e tan f ln tan þ
2 4
where N, corresponding to the mean anomaly M of elliptic orbits, is defined as

follows
" #12
l
N¼ 3
t tp
ðaÞ
The formulae given above use can be expressed in terms of hyperbolic (instead
of circular) functions, as follows. Remembering the definitions
ex þ ex ex ex
cosh x ¼ sinh x ¼
2 2
it is easy to verify that the following equality
cosh2 x sinh2 x ¼ 1
is an identity. Thus, the parametric equations of a hyperbola, with respect to a

parameter H, can be written as follows
x ¼ a cosh H
y ¼ b sinh H
where a and b are the transverse semi-axis and the conjugate semi-axis of the
hyperbola. Likewise, the radius vector from the occupied focus F to a point P of a
hyperbolic trajectory is expressible as follows
r ¼ að1 e cosh H Þ
The identities shown above involving the true anomaly / and the Gudermannian
f are also expressible in terms of / and H, by substituting
tan f ¼ sinh H
sec f ¼ cosh H
into the expressions shown above

1
e þ cos / ðe2 1Þ2 sin/
sec f ¼ tan f ¼
In addition, since

f H
tan ¼ tanh
2 2
then the expression

1
/ eþ1 2 f
tan ¼ tan
2 e1 2
becomes
1
/ eþ1 2 H
tan ¼ tanh
2 e1 2
The definition given above
ex þ ex ex ex
cosh x ¼ sinh x ¼
2 2
solved for ex yields
ex ¼ cosh x þ sinh x
which, in turn, taking the natural logarithms of both sides, yields
x ¼ lnðcosh x þ sinh xÞ
Now, since (as has been shown above)
tan f ¼ sinh H
sec f ¼ cosh H
then
H ¼ lnðtan f þ sec fÞ
From the following trigonometric identities
2t 1 þ t2 tan a þ tan b
tan a ¼ sec a ¼ tanða þ bÞ ¼
1 t2 1 t2 1 tan a tan b
where t = tan(a/2), it follows that

2t 1 þ t2
H ¼ lnðtan f þ sec fÞ ¼ ln þ
1 t2 1 t2
2 p 3
f
tan þ tan
1þt 6 4 2 7 f p
¼ ln ¼ ln6
4 p 7 ¼ ln tan þ
1t f 5 2 4
1 tan tan
4 2
Consequently, the relation between time and H, in case of a hyperbolic trajec-

tory, is
N ¼ e sinh H H
By comparing this relation with Kepler’s equation for elliptic orbits shown in the
preceding section, that is,
M ¼ Æ e sin Æ
setting
Æ ¼ iH
M ¼ iN
where i = (−1)½, and remembering the identities
sinðxÞ ¼ sin x
sinðixÞ ¼ i sinh x
there results
N ¼ e sinh H H
The function

f p
H ¼ ln tan þ
2 4
expresses H as a function of f. The inverse function, which expresses f as a

function of H, is called the Gudermannian of H and is obtained by solving the
preceding expression for f as follows

f p f p p
e ¼ tan
H
þ arctan eH ¼ þ f ¼ 2 arctan eH
2 4 2 4 2
As has been done with the elliptic orbits, we consider now the hyperbolic
trajectory actually traced by a spacecraft.
With reference to the preceding figure, let us consider a given hyperbolic tra-
jectory (of equation x*2/a2 − y*2/b2 = 1) in a reference system x*y* having its
origin O in the centre of the hyperbola. In the same reference system, let us consider
an auxiliary equilateral hyperbola (of equation x*2/a2 − y*2/a2 = 1), such that its
centre coincides with the centre O of the hyperbolic trajectory and its asymptotes
are orthogonal to each other and form angles of 45° with respect to the line of
apsides (x*-axis).
Let P be the position occupied, at a given time t, by a spacecraft along its
hyperbolic trajectory. Let r and / be the polar co-ordinates of P with respect to a
reference system having its origin in the occupied focus F of the hyperbolic tra-
jectory. Let us produce a straight line QR passing through P and perpendicular to
the line of apsides, Q and R being the points where this straight line intersects,
respectively, the auxiliary equilateral hyperbola and the line of apsides. Let V be the
vertex of the hyperbolic trajectory. Let S be the (shaded) area of the sector OQV
bounded by the rectilinear segments OV and OQ and by the arc QV of the auxiliary
equilateral hyperbola.
Because of the definition given above for the quantity H, this area is equal to
1
S ¼ AreaðOQVÞ ¼ a2 H
2
In addition, since the Cartesian co-ordinates of Q are x* OR, y* RQ, and

the auxiliary equilateral hyperbola satisfies the condition x*2/a2 − y*2/a2 = 1, then
there results
OR ¼ a cosh H
RQ ¼ a sinh H
The quantity H (which corresponds to the eccentric anomaly Æ of an elliptic

orbit) is not an angle. It is a dimensionless real number, proportional to the area S of
the hyperbolic sector defined above. H is positive or negative depending on whether
the hyperbolic sector OQV is placed above or below the x*-axis.
In case of a parabolic trajectory, Sect. 1.2 has shown Barker’s equation

/ /
tan3 þ 3 tan ¼ 2B
2 2
where
12
l
B¼3 3 t tp
p
By defining the quantity

1 /
D ¼ p2 tan
2
where D corresponds to the eccentric anomaly in case of parabolic trajectories,

Barker’s equation can be written in the following form
3
1 D
t tp ¼ 1 þ pD
2l2 3
and then the interval of time tf − ti necessary for a spacecraft to move along a
parabolic trajectory from an initial point Pi to a final point Pf is
" ! #
1 D3f D3i
tf ti ¼ 1 þ pDf þ pDi
2l2 3 3
1.5 The Lambert Problem
Section 1.3 has shown how to compute the time taken by a given planet revolving
about the Sun (or by an artificial satellite revolving about the Earth) to move along
its orbit from a given initial position to another given final position.
The present section shows how to solve the Lambert problem, which consists in
determining the transfer trajectory which a spacecraft must follow to move, under
the action of impulsive thrust manoeuvres and the gravitational force exerted by the
Sun, from a given planet (e.g., the Earth) to another given planet (e.g., Mars) within
a given interval of time tf − ti and in a given direction of flight. With reference to
the following figure, due to the courtesy of NASA [13], we consider a planetary
transfer from the Earth to Mars.
Given two points Pi and Pf in space, the time of transfer tf − ti and the direction
of motion, we want to determine the conic section which a spacecraft must follow
to move from Pi to Pf in the given interval of time tf − ti. The same results are also
applicable to the transfer trajectory of a spacecraft going from one to another orbit
around the Earth; such is the case with an artificial satellite to be injected from a
low-Earth orbit into a higher-altitude (e.g., geostationary) orbit.
Let ri and rf be the magnitudes of the position vectors of a spacecraft with respect
to a heliocentric reference system at the moment of, respectively, its departure from
the point Pi and its arrival at the point Pf. Let D/ be the angle through which the
position vector of the spacecraft along the transfer trajectory turns from the initial
point to the final point. Lambert’s theorem states [14] that the transfer time of a
body moving between two points of a conic trajectory depends only on the sum of
the distances of the two points from the origin of force, the linear distance between
such points, and the major semi-axis of the conic trajectory, that is,

Dt ¼ f ri þ rf ; c; a
1.5 The Lambert Problem 39
where c and a designate, respectively, the chord (or linear distance between the two
points) and the major semi-axis of the transfer trajectory, as shown in the following
figure, where the Sun is assumed to be the origin of force.
This figure illustrates two different cases for the same occupied focus F, the same
(counterclockwise) direction of motion, and the same a, ri, rf, c, and D/. There are
two possible different elliptic transfer trajectories between the departure and the
arrival points, with two associated different transfer times. This is because the
vacant focus F′ is in different positions in the two cases.
In the first case, shown on the left-hand side of the preceding figure, the vacant
focus F′ falls outside the plane surface (the shaded area) delimited by the chord
c and the transfer conic trajectory; in the second case, shown on the right-hand side,
the vacant focus F′ falls inside the shaded area.
Let Pi (of polar co-ordinates ri, /i) and Pf (of polar co-ordinates rf, /f) be,
respectively, the initial point and the final point of the transfer conic trajectory. The
length c of the chord joining Pi with Pf depends on ri, rf, and the heliocentric
transfer angle D/ /f − /i, through which the position vector r turns between Pi
and Pf, according to the well-known law of cosines:
c2 ¼ ri2 þ rf2 2ri rf cosðD/Þ

Now, ri, rf, and c are known from a preliminary mission analysis. If the major
semi-axis a of the transfer conic trajectory were determined, then the transfer time
tf − ti would also be determined, as will be shown below.
Following again Jordan [14], when the transfer trajectory is an ellipse travelled
counterclockwise, with the vacant focus F′ falling outside the shaded area, then the
transfer time Dt is given by
3 12
a
Dt ¼ ½ða sin aÞ ðb sin bÞ
l
where a is the major semi-axis of the elliptic transfer trajectory having the Sun at
the occupied focus F, l is the gravitational parameter of the Sun, a and b (such that
0 a p and 0 b p) are, respectively, first and second Lambert’s angles
for elliptic transfer trajectories, defined by cos a = 1 − s/a, cos b = 1 − (s − c)/a,
and s = (ri + rf + c)/2 is the semi-perimeter of the triangle whose sides are ri, rf,
and c. When the transfer trajectory is an ellipse travelled counterclockwise, with the
vacant focus F′ falling inside the shaded area, the transfer time Dt is given by
3 12
a
Dt ¼ T ½ða sin aÞ ðb sin bÞ
l
where T = 2p(a3/l)1/2 is the period of the elliptic transfer orbit.

When the transfer trajectory is a parabola, then the transfer time is
1
1 2 2h 3 3
i
Dt ¼ s 2
ð s cÞ 2
3 l
where, as shown above, s is the semi-perimeter of the triangle whose sides are ri, rf,
and c. The minus-or-plus sign in the preceding expression is due to the fact that, for
given values of ri, rf, and D/, there are two parabolas having their focus in the Sun
and joining the points Pi and Pf. The upper sign (minus) takes effect for the parabola
which satisfies the condition D/ < p (such that F is in the unshaded area), as
follows
1
1 2 2h 3 3
i
Dt ¼ s 2 ð s cÞ 2
3 l
Likewise, the lower sign (plus) takes effect for the other parabola, which satisfies
the condition D/ > p (such that F is in the shaded area), as follows
1
1 2 2h 3 3
i
Dt ¼ s 2 þ ð s cÞ 2
3 l
When D/ = p, the transfer time is

1
1 2 2 3
Dt ¼ s2
3 l
Likewise, when the transfer trajectory is a hyperbola, there are two hyperbolas
having one (F) of their foci in the Sun and joining the points Pi and Pf. For the
hyperbola such that the chord-flight path area does not shade F (D/ < p), the time
of transfer is
3 12
a
Dt ¼ ½ðsinh a aÞ ðsinh b bÞ
l
where cosh a = 1 − s/a and cosh b = 1 − (s − c)/a.

For the hyperbola such that the chord-flight path area does shade F (D/ > p), the
time of transfer is
3 12
a
Dt ¼ ½ðsinh a aÞ þ ðsinh b bÞ
l
Finally, when D/ = p, the two possible hyperbolas are identical, so that s = 0

and b = 0, and the time of transfer is
3 12
a
Dt ¼ ðsinh a aÞ
l
The type (ellipse, parabola, or hyperbola) of the conic transfer trajectory depends
on the velocity vi of the spacecraft at the departure time ti. Remembering the
expression (shown in Sect. 1.1) of the mechanical energy per unit mass possessed
by the spacecraft
1 l
E ¼ v2
2 r
where r is the radius vector of the spacecraft with respect to the Sun, l is the
gravitational parameter of the Sun and v the velocity of the spacecraft with respect
to the Sun, the type of conic section followed by the spacecraft during its transfer
from a planet to another depends on whether
E\0 ðelliptic transfer trajectoryÞ

E ¼ 0 ðparabolic transfer trajectoryÞ
E [ 0 ðhyperbolic transfer trajectoryÞ
The condition given above, in terms of the velocity vi of the spacecraft with
respect to the Sun at the departure time, becomes
1
vi \ð2l=r Þ =2 ðelliptic transfer trajectoryÞ
1
vi ¼ ð2l=r Þ =2 ðparabolic transfer trajectoryÞ
1
vi [ ð2l=r Þ =2 ðhyperbolic transfer trajectoryÞ
Now, as has been shown above, vi, ri, rf, c, D/ /f − /i, and Dt are known
from a preliminary mission analysis. The problem is the determination of the major
semi-axis, or transverse semi-axis, of a chosen transfer conic trajectory (ellipse or
parabola or hyperbola) corresponding to the known data. To this end, it is necessary
to solve iteratively an algebraic equation of the type f(a) = 0, using one of the
methods which will be shown below.
First of all, it is necessary to establish the limits of the range within which the
unknown quantity a can vary. This operation, called bracketing, is aimed at finding
two abscissae a2 and a1 (with a2 < a1) such that f(a2) and f(a1) should have opposite
signs, that is, such that the value of the product f(a1)f(a2) should be less than zero.
When this condition is satisfied, then we are assured of the existence of least one
root, a, contained in the interval [a2, a1] such that f(a) = 0.
The method of bisection is a sequence of operations aimed at halving the interval
of search. At the start of the sequence, let a[1] [1]
2 and a1 be, respectively, the lower
endpoint and upper endpoint of the initial interval a2 a a1, such that the
condition f(a[1] [1]
1 )f(a2 ) < 0 is satisfied. Let us estimate the unknown root as the point
which bisects this interval, that is, let us take
½1 ½1
½1 a2 þ a1
a0 ¼
2
In case of f(a[1] [1]

0 ) and f(a2 ) having the same sign, the new interval of search
becomes [a2 , a1 ] = [a0 , a[1]
[2] [2] [1] [2] [2] [1] [1]
1 ]; otherwise, it becomes [a2 , a1 ] = [a2 , a0 ].
Now, the unknown root is estimated again as follows
½2 ½2
½2 a2 þ a1
a0 ¼
2
The sequence of operations indicated above is repeated n times, until the

1 − a2 converges to zero within the desired
absolute value of the difference a[n] [n]
accuracy, which always happens if the given function f(a) is continuous. This
method is linearly convergent, the number of the correct decimal places in the
computed value being proportional to the number, n, of iterations.
The tangent (or Newton–Raphson) method requires the knowledge of the first
derivative df/da of the function f f(a) to compute a sequence of points a[1],
a[2], … converging to the unknown root a. The sequence of operations is
indicated below. First, take a guess a[1] of the unknown value a. Then,
approximate f(a) around a[1] by means of a Taylor-series expansion truncated
after its second term, that is, in geometrical terms, f(a) is approximated around
a[1] by means of the straight line which passes through the point of co-ordinates
a[1], f(a[1]) and has the same slope as the curve f(a), as follows

½1 df
f ðaÞ f a þ a a½1 þ
da 1
Now, a better approximation, a[2], than a[1] is obtained by equating this truncated
expansion to zero, which yields

½1 df
f a þ a½2 a½1 ¼ 0
da 1
that is,

½2 ½1 f a½1
a ¼a
df
da 1
Now the absolute value of f(a[2]) is computed and compared with a chosen
tolerance e. If |f(a[2])| e, then the iterative process comes to an end, because the
approximation reached is judged satisfactory; otherwise, another approximate
value, a[k+1], is computed by means of the recursion

½k þ 1 f a½k
½k
a ¼a
df
da k
while |f(a[k])| > e. This method has, over the bisection method, the advantage of
having a quadratic, instead of linear, convergence. However, unlike the bisection
method, the Newton–Raphson method does not always converge, because it is
misled by local extrema (i.e., points of maximum or minimum) which may exist in
the interval of search. In such points, as is well known, the first derivative vanishes,
causing the value of the next approximation a[k+1] to go to infinity. This problem is
avoided if Newton–Raphson steps are alternated with bisection steps every time a
value computed by means of the Newton–Raphson method falls outside the interval
of search; this alternation of steps is not necessary if the function f(a) be known in
advance to have no local extrema in the interval where the root is sought.
In addition to the Newton–Raphson method, a method which does not require
the knowledge of the derivative of the function f(a) is the secant method, which
comprises the following steps:
• take a first estimate a[1] of the unknown root and compute the correspondent
value f(a[1]);
• take a second estimate a[2] of the unknown root and compute the correspondent
value f(a[2]);
• approximate the derivative df/da computed at a[2] by the slope of the straight
line which joins P1, whose co-ordinates are a[1], f(a[1]), with P2, whose
co-ordinates are a[2], f(a[2]); this means that the derivative df/da at a[2] is
approximated as follows

df f a½2 f a½1

da 2 a½2 a½1
• use (df/da)2, approximated as shown above, to compute

½3 ½2 f a½2
a ¼a
df
da 2
that is,
a½2 a½1
a½3 ¼ a½2 f a½2
f ða½2 Þ f ða½1 Þ
and compute f(a[3]);

• if |f(a[3])| e, then the iterative process terminates, because the approximation
reached is judged satisfactory; otherwise, another approximate value
a½3 a½2
a½4 ¼ a½3 f a½3
f ða½3 Þ f ða½2 Þ
is computed, and so on, k times, while the condition |f(a[k])| e is not satisfied.
The secant method converges more slowly than the tangent method and also
suffers from the fact that the denominator [f(a[k+1]) − f(a[k])] of the recurrence
a½k þ 1 a½k
a½k þ 2 ¼ a½k þ 1 f a½k þ 1
f ða½k þ 1 Þ f ða½k Þ
may go to zero, in which case the same measures shown for the Newton–Raphson
method must be taken.
Still another possible method to find the required value of a is Müller’s method
of parabolic interpolation, which uses a quadratic interpolation and consequently
requires three points in the vicinity of the root to be found. This method is briefly
described below.
Let a2 and a1 be the endpoints of an interval a2 a a1 containing the
unknown value, a, of the root. Let f(a) be the function which must vanish in the
point a. As shown above, the existence of at least one real root of f(a) = 0 lying
between a2 and a1 is assured if f2 f(a2) and f1 f(a1) have opposite signs.
Following Gerald and Wheatley [12], the computation is performed as follows:
• take a third point a0 laying between a2 and a1 and compute f0 f(a0);
• set h1 = a1 − a0, h2 = a0 − a2, and c = h2/h1;
• compute the coefficients
cf1 f0 ð1 þ cÞ þ f2
A¼
ch21 ð1 þ cÞ
f1 f0 Ah21
B¼
h1
C ¼ f0
of the interpolating parabola f(a) = A(a − a0)2 + B(a − a0) + C;

• compute the estimated root of f(a) = 0 as follows
2C
a ¼ a0 1
B ðB2 4AC Þ2
where the sign plus or minus is chosen so that the denominator should have the
maximum absolute value (i.e., if B > 0, choose plus; if B < 0, choose minus; if
B = 0, choose either) and compute f f(a);
• check the computed value of a to determine which set of three points should be
used in the next iteration (if a is greater than a0, take a0, a1, and the root a for the
next iteration; if a is less than a0, take a0, a2, and the root a for the next
iteration); and
• reset the subscripts 0, 1, and 2 so that a0 should be placed between a2 and a1,
and repeat the cycle while the computed value of a does not satisfy the condition
f(a) = 0 to some acceptable degree of tolerance.
Gedeon [15, 16] applies the Newton–Raphson method to compute iteratively the
major (or transverse) semi-axis of the transfer conic trajectory. To this end, he
introduces a new variable
s
z¼
2a
where s = (ri + rf + c)/2, and two constants, which are

c12 l 12
w¼ 1 ns ¼
s s3
Gedeon substitutes these quantities into the two following equations

3 12
a
Dt ¼ ½ð1 kÞmp þ k ða sin aÞ
ðb sin bÞ
l
(which holds for an elliptic transfer orbit travelled counterclockwise), where m is

the number of revolutions in the transfer orbit, and k = ± 1 indicates whether the
vacant focus F′ falls outside (upper sign) or inside (lower sign) the area enclosed by
the chord and the trajectory; and
3 12
a
Dt ¼ ½ðsinh a aÞ ðsinh b bÞ
l
(which holds for a hyperbolic transfer orbit with D/ < p).

By so doing, Gedeon obtains the following general expressions [16]:
!
1 1k h 1 1 1
i
N ¼ ns Dt ¼ 1 1 mp þ k f jzj2 jzj2 ð1 zÞ2
zjzj2 22 2
h 1
1 1 io
f wjzj2 wjzj2 1 w2 z 2
!" #
dN 1 k w3 3N
¼ 1 1
jzj22 ð1 zÞ2 ð1 w2 zÞ2 22
1 1
dz
where, in case of hyperbolic transfer trajectories, the function f(…) is argsinh and
k = + 1, whereas, in case of elliptic transfer trajectories, f(…) = arcsin and
k = ± 1.
In the latter case, k = + 1 when the vacant focus is outside the area enclosed by
the chord and the trajectory, and k = −1 when the vacant focus is inside this area.
The quantity w is positive when the heliocentric transfer angle D/ /f − /i is less
than p, and negative when D/ > p. The expressions given above are functions of
the unknown quantity z = s/(2a).
It is to be noted that the endpoints of an interval within which the unknown value
of the major semi-axis, or of z = s/(2a), is sought cannot be arbitrary, since there are
values of a for which the Lambert equations have no solutions. For example, there
is a minimum value (amin) of the major semi-axis such that, if a < amin, no transfer
ellipse is possible. This minimum value is given by
s
amin ¼
2
where s is the semi-perimeter of the triangle FPiPf.
1.6 Transfer Times for Elliptic, Parabolic, and Hyperbolic

Trajectories
Let Æi and Æf be the eccentric anomalies of, respectively, the initial (or departure)
point Pi and the final (or arrival) point Pf of an elliptic transfer orbit. Remembering
the equation for the time of transfer
1.6 Transfer Times for Elliptic, Parabolic, and Hyperbolic Trajectories 47
3 12
a

Dt ¼ 2mp þ Æ f e sin Æ f ðÆ i e sin Æ i Þ
l
derived in Sect. 1.3, we set m = 0, thereby excluding the possibility of multiple

revolutions of our spacecraft in its transfer ellipse having the Sun at one focus. For
m = 0, the previous expression becomes
3 12
a

Dt ¼ Æ f e sin Æ f ðÆ i e sin Æ i Þ
l
that is,
3 12
a

Dt ¼ Æ f Æ i e sin Æ f sin Æ i
l
Now, since the following expression

Æf Æi Æf þ Æi
sin Æ f sin Æ i ¼ 2 sin cos
2 2
is a trigonometric identity, then the expression for the time of transfer becomes
3 12
a Æf Æi Æf Æi Æf þ Æi
Dt ¼ 2 e sin cos
l 2 2 2
Since the eccentric anomalies appear in the previous expression only as their
semi-sum and semi-difference, we can define the two following quantities
Æf Æi
w¼
2

Æf þ Æi
cos v ¼ e cos
2
We have then
3 12
a
Dt ¼ 2 ðw sin w cos vÞ
l
The other quantities (ri, rf, and c) involved in Lambert’s theorem can also be
expressed as a function of w and v, as will be shown below.
Following Battin [11], the equation of orbit r = a(1 − e cos Æ) makes it pos-
sible to express the sum of the two radii, ri and rf, as follows

ri þ rf ¼ að1 e cos Æ i Þ þ a 1 e cos Æ f ¼ 2a ae cos Æ f þ cos Æ i
Now, since

Æf þ Æi Æf Æi
cos Æ f þ cos Æ i ¼ 2 cos cos
2 2

Æf þ Æi Æf Æi
ri þ rf ¼ 2a 1 e cos cos ¼ 2að1 cos v cos wÞ
2 2
The following expressions of Sect. 1.3

/ að 1 þ e Þ 2 Æ
sin2
¼ sin
2 r 2

/ að 1 e Þ 2 Æ
cos2 ¼ cos
2 r 2
can also be written as follows

1 / 1 Æ 1 / 1 Æ
r sin
2 ¼ ½að1 þ eÞ sin
2 r cos
2 ¼ ½að1 eÞ cos
2
2 2 2 2
These expressions, written for ri and rf, are

1 /i 1 Æi 1 /i 1 Æi
ri2 sin ¼ ½að1 þ eÞ2 sin ri2 cos ¼ ½að1 eÞ2 cos
2 2 2 2

1 /f 1 Æf 1 / f 1 Æf
rf2 sin ¼ ½að1 þ eÞ2 sin rf2 cos ¼ ½að1 eÞ2 cos
2 2 2 2
Therefore,

1 /f /i 1 /f / /f /
ri rf 2 cos ¼ ri rf 2 cos cos i þ sin sin i
2 2 2 2 2

Æf Æi Æf Æi
¼ a ð1 eÞcos cos þ ð1 þ eÞsin sin
2 2 2 2

Æf Æi Æf þ Æi
¼ a cos e cos ¼ aðcos w cos vÞ
2 2
This is because we have set
Æf Æi
w¼
2
Æf þ Æi
cos v ¼ e cos
2
In addition, we note that

ri þ rf ¼ að1 e cos Æ i Þ þ a 1 e cos Æ f ¼ 2að1 cos v cos wÞ
This makes it possible to express the chord length, c, as a function of w and v, as

follows
2
c2 ¼ ri2 þ rf2 2ri rf cosðD/Þ ¼ ri þ rf 2ri rf 2ri rf cosðD/Þ

2 D/
¼ ½2að1 cos v cos wÞ 2ri rf 2ri rf cos 2
2

D/
¼ 4a2 1 2 cos v cos w þ cos2 v cos2 w 2ri rf 2ri rf cos2
2

D/
sin2 ¼ 4a2 1 2 cos v cos w þ cos2 v cos2 w 2ri rf
2

D/ D/
1 þ cos2 sin2 ¼ 4a2 1 2 cos v cos w þ cos2 v cos2 w
2 2

D/
2ri rf 2cos2 ¼ 4a2 1 2 cos v cos w þ cos2 v cos2 w
2
" #
a ðcos w cos vÞ2
2
4ri rf ¼ 4a2 ð1 2 cos v cos w þ cos2 v cos2 w
ri rf

cos2 w þ 2 cos v cos w cos2 vÞ ¼ 4a2 sin2 w 1 cos2 w cos2 v

¼ 4a2 sin2 w sin2 w cos2 v ¼ 4a2 sin2 w sin2 v
Therefore,
c ¼ 2a sin w sin v
Since w and v can be expressed as functions of a, ri + rf, and c by means of the

formulae given above, then these functions of a, ri + rf, and c can be substituted
into the equation
3 12
a
l
This proves Lambert’s theorem, according to which the time of transfer, Dt,
depends only on a, ri + rf, and c.
Lagrange’s form of the equation for the transfer time results from introducing
two functions of w and v, defined as follows
a ¼ vþw
b¼vw
that is,
ab
w¼
2
aþb
v¼
2
so that
2s ¼ ri þ rf þ c ¼ 2að1 cos v cos w þ sin v sin wÞ ¼ 2a½1 cosðv þ wÞ

¼ 2að1 cos aÞ
2ðs cÞ ¼ ri þ rf c ¼ 2að1 cos v cos w sin v sin wÞ
¼ 2a½1 cosðv wÞ ¼ 2að1 cos bÞ
where s is the semi-perimeter of the triangle FPiPf. It follows that

s a s
cos a ¼ 1 sin2 ¼
a 2
2a
sc b sc
cos b ¼ 1 sin2 ¼
a 2 2a
The equation expressing the time of transfer, that is,

3 12
a
l
can also be written in terms of a and b as follows

3 12 3 12
a a ab ab aþb
Dt ¼ 2 ðw sin w cos vÞ ¼ 2 sin cos
l l 2 2 2
3 12 3 12
a ab aþb a
¼ a b 2sin cos ¼ ða b sin a þ sin bÞ
l 2 2 l
3 12
a
¼ ½ða sin aÞ ðb sin bÞ
l
which is the same expression shown in the previous section, in case of the transfer
trajectory being an ellipse travelled counter clockwise, with the vacant focus F′
falling outside the shaded area.
The previous expression takes account of the following trigonometric identity

ab aþb
sin a sin b ¼ 2sin cos
2 2
When the transfer trajectory is a parabola, then the time of transfer results from
Barker’s equation derived in Sect. 1.2, that is,

/ /
tan3 þ 3tan ¼ 2B
2 2
Remembering that B = 3(l/p3)½(t − tp) and setting r = p½ tan(//2), Barker’s

equation can be written as follows
1
6l2 t tp ¼ r3 þ 3pr
The preceding expression, written for, respectively, the final point, Pf, and the
initial point, Pi, of a parabolic transfer obit, is
1 1
6l2 t tp ¼ r3f þ 3prf 6l2 t tp ¼ r3i þ 3pri
Subtracting the second equation from the first yields

1
6l2 tf ti ¼ r3f r3i þ 3p rf ri
Since (a3 − b3) = (a − b)(a2 + ab + b2) is an algebraic identity, then

1
6l2 tf ti ¼ rf ri r2f þ rf ri þ r2i þ 3p
In addition, since r2f + r2i = (rf − ri)2 + 2rf ri is also an algebraic identity, then
1 h 2 i
6l2 tf ti ¼ rf ri rf ri þ 3 p þ rf ri
Now, remembering the definition given above

1 /
r ¼ p tan
2
2
there results

1 /f 1 /i /f /
p þ rf ri ¼ p þ p tan
2 p tan
2 ¼ p 1 þ tan tan i
2 2 2 2
Remembering the following trigonometric identities
1
sin a sin b ¼ ½cosða bÞ cosða þ bÞ
2
1
cos a cos b ¼ ½cosða bÞ þ cosða þ bÞ
2
and dividing the first by the second, there results
cosða bÞ cosða þ bÞ
tan a tan b ¼
cosða bÞ þ cosða þ bÞ
In the present case, there results

/ /i /f þ /i
cos f cos
/f / 2 2
tan tan i ¼
2 2 / f /i /f þ /i
cos þ cos
2 2
then

/f /i
2cos
/f / 2
1 þ tan tan i ¼
2 2 /f /i /f þ /i
cos þ cos
2 2
Now, remembering the trigonometric identity

ab aþb
cos a þ cos b ¼ 2 cos cos
2 2
there results

/f /i /f þ /i /f /
cos þ cos ¼ 2 cos cos i
2 2 2 2
hence
2 32 3
/f /i
2 cos
/f / 6 1 76 7
p 1 þ tan tan i ¼p
2
¼ p6 7 6 1 7
2 2 /f / 4 / 5 4 /f 5
2 cos cos i cos i cos
2 2 2 2

/f /i / /f /f /i
cos ¼ p sec i sec cos
2 2 2 2
that is,

/i /f /f /i
p þ ri rf ¼ p sec sec cos
2 2 2
Remembering now the equation of a parabolic trajectory (e = 1)
p p 1
r¼ ¼ ¼p
1 þ cos / / / / / /
cos2 þ sin2 þ cos2 sin2 2 cos2
2 2 2 2 2

1 /
¼ p sec2
2 2
and solving the preceding equation for sec(//2), there results

12 12
/ 2 12 /f 2 12
sec i ¼ r sec ¼ r
2 p i 2 p f
These expressions, in turn, inserted into

/i /f /f /i
p þ ri rf ¼ p sec sec cos
2 2 2
yield after simplification

1 /f /i
p þ ri rf ¼ 2 ri rf 2 cos
2
From the preceding figure and the law of cosines
c2 ¼ ri2 þ rf2 2ri rf cosðD/Þ
there results

D/ 2 D/
c ¼
2
ri2 þ rf2
2ri rf cosðD/Þ ¼ ri2 þ rf2
2ri rf cos 2
sin
2 2

2 D/ 2 D/
¼ ri þ rf 2ri rf cos
2 2
1 þ cos ¼ ri2 þ rf2 2ri rf
2 2

2 D/ 2 D/
2 cos 1 ¼ ri þ rf 4ri rf cos
2 2
þ 2ri rf
2 2

2 D/
¼ ri þ rf 4ri rf cos2
2
Now, remembering the definition of semi-perimeter
ri þ rf þ c
s¼
2
there results
2s ¼ ri þ rf þ c
ri þ rf ¼ 2s c
hence

2 D/ D/
c2 ¼ ri þ rf 4ri rf cos2 ¼ ð2s cÞ2 4ri rf cos2
2 2

D/
4ri rf cos2
¼ ð2s cÞ2 c2 ¼ 4s2 4sc þ c2 c2 ¼ 4sðs cÞ
2
Taking the square root of the expression written above yields

12 D/ 1
2 ri rf cos ¼ 2½sðs cÞ2
2
hence

1 /f /i 1
p þ ri rf ¼ 2 ri rf 2 cos ¼ 2½sðs cÞ2
2
where the upper sign (plus) or the lower sign (minus) takes effect when the
heliocentric transfer angle D/ = /f − /i is, respectively, less than or greater than p
radians.
So far, we have expressed the term p + ri rf appearing in the time equation
1 h 2 i
6l2 tf ti ¼ rf ri rf ri þ 3 p þ rf ri
for a parabolic transfer trajectory as a function of the variables ri + rf and c.

Now the term rf − ri appearing in this equation will be expressed as a function
of the same variables. To this end, we observe that the equation of a parabolic
transfer trajectory

p 1 /
r¼ ¼ p sec2
1 þ cos / 2 2
may also be written as follows
p þ r2
r¼
2
where r = p1/2 tan(//2). This is because

2 3
/ /
sin2 1 cos2
/ 2 2 6 1 7
r2 ¼ p tan2 ¼p ¼p ¼ p6
4 17
5
2 / / /
cos2 cos2 cos2
2 2 2

/
¼ p sec2 1
2
hence

/ r2
sec2
¼ þ1
2 p
this, in turn, introduced into the parabolic trajectory equation

1 2 /
r ¼ p sec
2 2
yields
p þ r2
r¼
2
It follows that

2 ri þ rf ¼ p þ r2i þ p þ r2f ¼ 2p þ r2i þ r2f ¼ 2p þ 2ri rf 2ri rf þ r2i þ r2f
2 2
¼ 2p þ 2ri rf þ rf ri ¼ rf ri 2 p þ ri rf
which, solved for (rf − ri)2, yields

2

rf ri ¼ 2 ri þ rf p þ ri rf
Now, remembering that
ri þ rf ¼ 2s c
1
p þ ri rf ¼ 2½sðs cÞ2
there results
2 n 1
o
rf ri ¼ 2 2s c
2½sðs cÞ2
Since
h1 i
1 2 1 1
s2
ðs cÞ2 ¼ s þ ðs cÞ
2½sðs cÞ2 ¼ 2s c
2½sðs cÞ2
then
2 h1 i
1 2
rf ri ¼ 2 s2
ðs cÞ2
In addition, since rf − ri is never less than zero, then

1
h1 1
i
rf ri ¼ 22 s2
ðs cÞ2
Now, rf − ri has also been expressed as a function of the variables ri + rf, and c.
1 1 1 1
Finally, by introducing p þ ri rf ¼ 2½sðs cÞ2 and rf ri ¼ 22 ½s2
ðs cÞ2
1 h 2 i
into the time equation 6l2 tf ti ¼ ðrf ri Þ rf ri þ 3ðp þ ri rf Þ , there
results
1 1
h1 1
i h 1 i
1 2 1

6l2 tf ti ¼ 22 s2
ðs cÞ2 2 s2
ðs cÞ2 6½sðs cÞ2
Since the semi-perimeter s is defined as follows
ri þ rf þ c
s¼
2
then there results

12
1 ri þ rf þ c u
s ¼
2
1 ¼ 1
2 2 22
12
1 ri þ rf c v
ð s cÞ 2 ¼ 1 ¼ 1
22 22
where the following variables

1
u ¼ ri þ rf þ c 2
1
v ¼ ri þ rf c 2
are used for convenience. With these variables, the time equation
1 1
h1 1
i h 1 i
1 2 1

6l2 tf ti ¼ 22 s2
ðs cÞ2 2 s2
ðs cÞ2 6½sðs cÞ2
becomes
( )
1 1 u v u v 2 u v
6l tf ti ¼ 22 1
1 2 1
1 6 1
2
1
22 22 22 22 22 22
h i
¼ ðu
vÞ ðu
vÞ2 3uv ¼ u3
v3
that is, using again the variables ri + rf and c instead of u and v,

1 3 3
6l2 tf ti ¼ ri þ rf þ c 2
ri þ rf c 2
The time equation, as a function of semi-perimeter and chord, is

1
1 2 2h 3 3
i
tf ti ¼ s 2
ð s cÞ 2
3 l
which accounts for the expression given in Sect. 1.5. The time equation for a
parabolic transfer trajectory was found by Euler in 1743.
1.7 A Unified Form of Lambert’s Equations
In the preceding sections, the three basic conic sections (ellipse, parabola and
hyperbola) have been considered separately. This requires the previous selection of
a specific type of conic section for the transfer trajectory from an initial point Pi to a
final point Pf. Otherwise, the problem may be posed in general terms as follows:
given the position (ri) and velocity (vi) vectors of a spacecraft at an initial time ti,
determine the position (rf) and velocity (vf) vectors of the same spacecraft at a final
time tf, knowing that the spacecraft is only subject to impulsive thrust manoeuvres
and to the constant gravitational force exerted by a central body (e.g., by the Sun),
but without the previous knowledge of the type of conic section that the spacecraft
will follow. In the particular case of a planetary transfer, we want to determine the
correct values of a, Æi, Æf, and e which correspond to a desired transfer time tf − ti.
By using the expression, derived in Sect. 1.3, for the time of transfer in an elliptic
trajectory, that is,
3 12
a

Dt tf ti ¼ 2mp þ Æ f e sin Æ f ðÆ i e sin Æ i Þ
l
and dropping the term 2 mp, relating to a number m of revolutions made by the
spacecraft in its transfer orbit, there results
3 12
a

Dt tf ti ¼ Æ f e sin Æ f ðÆ i e sin Æ i Þ
l
or, setting DÆ Æf − Æi,

1 3

l2 Dt ¼ a2 DÆ þ e sin Æ i sin Æ f
Now, following Bate et al. [10], Vallado [17], and Anderson [18], we introduce
the following variables
1.7 A Unified Form of Lambert’s Equations 59
1 1
x ¼ a2 Æ f Æ i a2 DÆ
DÆ sin DÆ
S¼
ðDÆ Þ3
This yields
3
x3 S ¼ a2 ðDÆ sin DÆ Þ
hence
3 3
x3 S þ a2 sin DÆ ¼ a2 DÆ
Substituting the preceding expression into

1 3

l2 Dt ¼ a2 DÆ þ e sin Æ i sin Æ f
yields
1 3 3
l2 Dt ¼ x3 S þ a2 sin DÆ þ a2 e sin Æ i sin Æ f
Since

sin DÆ sin Æ f Æ i ¼ sin Æ f cos Æ i cos Æ f sin Æ i
is an identity, then
1 3 3
l2 Dt ¼ x3 S þ a2 sin DÆ þ a2 e sin Æ i sin Æ f
becomes
1 3 3
l2 Dt ¼ x3 S þ a2 sin Æ f cos Æ i cos Æ f sin Æ i þ a2 e sin Æ i sin Æ f
that is,
1 3
l2 Dt ¼ x3 S þ a2 sin Æ f cos Æ i cos Æ f sin Æ i þ e sin Æ i e sin Æ f
or
1 3

l2 Dt ¼ x3 S þ a2 e cos Æ f sin Æ i ðe cos Æ i Þsin Æ f
By multiplying and dividing the term within square brackets by

1
1 e2 2 ð1 e cos Æ i Þ 1 e cos Æ f
it is possible to write this term as follows

(" 1 # " 1 # )
ð1 e2 Þ2 sin Æ i e cos Æ f ð1 e2 Þ2 sin Æ f e cos Æ i

1 e cos Æ i 1 e cos Æ f 1 e cos Æ f 1 e cos Æ i

ð1 e cos Æ i Þ 1 e cos Æ f
1
ð 1 e2 Þ 2
Now, remembering the following formulae of Sect. 1.3

1
cos Æ e ð1 e2 Þ2 sin Æ
cos / ¼ sin/ ¼
1 e cos Æ 1 e cos Æ
the term within square brackets becomes

" #

ðsin /i Þ cos /f sin /f ðcos /i Þ 1
ð 1 e2 Þ 2
that is,
" #

sin /f /i 1
ð 1 e2 Þ 2
Thus, the equation of time

1 3

l2 Dt ¼ x3 S þ a2 e cos Æ f sin Æ i ðe cos Æ i Þsin Æ f
becomes
" #
1 3
l Dt ¼ x S þ a sin /f /i
2 3 2
1
ð1 e2 Þ2
In addition, since
a2 ¼ a1 a1 a 2
3 1
then the equation of time becomes

( )
1
að1 e cos Æ i Þa 1 e cos Æ f
l Dt ¼ x S þ sin /f /i
2 3
1
½ að 1 e 2 Þ 2
Now, since r = a(1 − e cos Æ), then

( )
1
ri rf
l Dt ¼ x S þ sin /f /i
2 3
1
½ að 1 e 2 Þ 2
This in turn, multiplying and dividing the second term on the right-hand side by
[1 − cos(D/)]½, yields
( )( 1
)
1 ri rf sinðD/Þ ½1 cosðD/Þ2
l2 Dt ¼ x S þ
3
1 1
½ að 1 e 2 Þ 2 ½1 cosðD/Þ2
1 1
Since ri rf ¼ ri rf 2 ri rf 2 , then the expression written above becomes
( 1 )( 1 1)
1 ri rf 2 sinðD/Þ ri rf 2 ½1 cosðD/Þ2
l2 Dt ¼ x S þ
3
1 1
½1 cosðD/Þ2 ½ að 1 e 2 Þ 2
Now, introducing two variables defined as follows

1
ri rf 2 sinðD/Þ ri rf ½1 cosðD/Þ
A¼ y¼
½1 cosðD/Þ
1
2 að 1 e 2 Þ
the equation of time becomes

1 1
l2 Dt ¼ x3 S þ Ay2
which, solved for Dt, yields the time equation for an interplanetary transfer in terms
of the universal variables x, S, A, and y
1
Dt ¼ x3 S þ Ay2 l2
1
Since

1 1 1
sinðD/Þ ¼ 1 cos2 ðD/Þ 2 ¼ ½1 þ cosðD/Þ2 ½1 cosðD/Þ2
where the upper sign (plus) takes effect when 0 < D/ < p, and the lower sign
(minus) takes effect when p < D/ < 2p, then A can be expressed as a function of ri,
rf and cos(D/) only as follows
1
ri rf 2 sinðD/Þ 1
A¼ 1 ¼ ri rf ½1 þ cosðD/Þ 2
½1 cosðD/Þ 2
The universal variables can be expressed as a function of the customary variables

and of z (DÆ)2, as follows
y 12
x¼
C
1 þ 1 for D/\p
A ¼ k ri rf ½1 þ cosðD/Þ 2 where k ¼
1 forD/ [p
8
>
1
> 1 cos z 2
>
>
>
> for z [ 0
AðzS 1Þ < z
y ¼ ri þ rf þ where C ¼ 12 for z ¼ 0
1
C2 >
>
>
> cosh z
1
1
>
>
2
: for z\0
8 1 ðz Þ1
>
> z2 sin z2
>
>
>
> for z [ 0
>
<
3
z2
S ¼ 6 for z ¼ 0
1
>
>
> 1 1
> sinh z z
> 2 2
>
>
: 3 for z\0
ðzÞ2
The correct value of z must be found iteratively, until the desired time of transfer
is obtained. The iterations can be performed by using one of the methods (Newton–
Raphson, secant, …) described in Sect. 1.5. At each iteration, a new value is
computed for z, and hence the correspondent values of x, S, and y. The iterative
process terminates when we find the value of z which corresponds, within the limits
of a chosen tolerance, to the desired transfer time.
Bate et al. [10], Vallado [17], and Anderson [18] indicate the following algo-
rithm to compute iteratively z. Given the following input values:
ri (initial position vector), rf (final position vector), Dt0 (desired transfer time),
and k (direction of motion, such that k = + 1 for D/ < p, and k = −1 for D/ > p)
The quantities f, g, and g′ are three of the four Lagrangian coefficients f, g, f′, and
g′, which depend on time through y. After the correct value of y, corresponding to
the input data (ri, rf, Dt0, and direction of motion), has been determined, it is
possible to compute the initial and final velocity vectors, that is,
rf f ri g0 r f r i
r0i ¼ r0f ¼
g g
where r′i and r′f are the velocity vectors at, respectively, the initial point (Pi) and the
final point (Pf) of the transfer trajectory. Thus, the components of such vectors, with
respect to a reference system XYZ whose origin O is in the centre of force, are
Xf fXi g0 Xf Xi
Xi0 ¼ Xf0 ¼
g g
Yf fXi g 0 Y f Yi
Yi0 ¼ Yf0 ¼
g g
Zf fXi g0 Zf Zi
Zi0 ¼ Zf0 ¼
g g
Other authors, for example Sharaf et alii [19] use similar algorithms to determine
iteratively z. The principal difference between Sharaf’s algorithm and the one
shown above resides in the fact that Sharaf et alii use a procedure of their own to
compute the Stumpff functions C(z) and S(z).
After the initial (r0i ) and final (r0f ) velocity vectors have been computed, the
required velocity increment Dv can be determined, as will be shown below.
Following Baun and Papadopoulos [20], let us consider, for example, a spacecraft
in a transfer conic trajectory from the Earth (subscript E) to Mars (subscript M).
Let vE and vM be, respectively, the orbital velocity of the Earth at the time (ti) of
departure and the orbital velocity of Mars at the time (tf) of arrival. The hyperbolic
excess velocity on the transfer trajectory at t = ti results from a vector subtraction
vE1 ¼ r0i vE
Likewise, the hyperbolic excess velocity on the transfer trajectory at t = tf results

from
vM1 ¼ r0f vM
The velocity increment Dv is the difference between the spacecraft velocity

around the planet and the spacecraft velocity with respect to the same planet.
Let the spacecraft be in a low-altitude circular orbit around the Earth at the time
(ti) of departure. Let rE be the radius of this circular orbit with respect to the centre
of the Earth. The velocity (vLEO) possessed by the spacecraft along its circular orbit
around the Earth is given by the vis-viva integral
12
lE
vLEO ¼
rE
At the end of the heliocentric transfer trajectory, at the time (tf) of arrival, the
spacecraft is to be injected into a low-altitude circular orbit around Mars. Let rM be
the radius of this circular orbit with respect to the centre of Mars. The velocity
possessed by the spacecraft (vLMO) along its circular orbit around Mars is
12
lM
vLMO ¼
rM
Thus, the velocity increment (Dvi) to be given to the spacecraft at the time of its
departure from the low-altitude circular orbit around the Earth results from
Dvi ¼ vE1 vLEO
Likewise, the velocity increment (Dvf) to be given to the spacecraft at the time of
its arrival to the low-altitude circular orbit around Mars results from
Dvf ¼ vM1 vLMO
The total increment of velocity (Dvtotal) to be given to the spacecraft is then
Dvtotal ¼ Dvi þ Dvf
1.8 An Example of Solution of Lambert’s Problem Using

Universal Variables
Let us consider the following example, proposed by Curtis [21]. A spacecraft

revolving about the Earth performs an orbit transfer to move from the initial point Pi
having Cartesian co-ordinates
Xi ¼ 5000 km
Yi ¼ 10000 km
Yi ¼ 2100 km
(with respect to the geocentric-equatorial reference system XYZ) to the final point Pf
having Cartesian co-ordinates
Xf ¼ 14600 km
Yf ¼ 2500 km
Zf ¼ 7000 km
with respect to the same reference system. Let Dt0 = 3600 s be the time interval
chosen to perform the orbit transfer. Let the direction of motion be prograde (or
easterly, or counter clockwise). The gravitational parameter of the Earth, which is
the central body in the present case, is l = 3.986 105 km3/s2, hence,
l½ = 6.31348 102 km3/2/s. The magnitudes ri and rf of, respectively, the initial
position vector (ri) and the final position vector (rf) are
h 2
2 2 i12
ri ¼ 5 103
þ 10 103 þ 2:1 103 ¼ 11:37585 103 km
h 2 2 2 i12
rf ¼ 14:6 103 þ 2:5 103 þ 7 103 ¼ 16:38322 103 km
The cosine of the angle formed by the two position vectors ri and rf is
rr rf 5 14:6 þ 10 2:5 þ 2:1 7

¼ ¼ 0:17867
ri rf 11:37585 16:38322
Hence, the geocentric transfer angle is
D/ ¼ arccosð0:17867Þ ¼ 1:75043 radians ¼ 100 :292
Since there results 0 < D/ < p, then

1
A ¼ þ ri rf ½1 þ cosðD/Þ 2 ¼ þ ½11:37585 103 16:38322 103
1
ð1 0:17867Þ2 ¼ 1:23723 104 km
The lower endpoint (zL) and the upper endpoint (zU) of the interval
zL \z\zU
within which the unknown value of z is sought are taken initially as follows
zL ¼ 4p2 ¼ 39:4784
zU ¼ 4p2 ¼ 39:4784
However, since in the present case A is positive, then y must also be positive.
Consequently, all values of z leading to negative values of y must be discarded. This
fact narrows the interval of search from zL < z < zU to a smaller width, as will be
shown below.
First, to determine the sign of y at the upper endpoint (zU) of the interval, we
compute the values of the Stumpff functions C(z) and S(z) at z = zU, as follows
1

1 cos 39:42
CU ¼ C ð39:4Þ ¼ ¼ 4:95 107
39:4
1 1
39:42 sin 39:42
SU ¼ Sð39:4Þ ¼ 3 ¼ 0:031
39:42
The corresponding value of y is

1
yU ¼ ri þ rf þ AðzU SU 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103
CU2
1
þ 1:23723 104 ð39:4 0:031 1Þ 1 ¼ 3:9211363 10
6
ð4:95 107 Þ2
Thus, at z = zU, there results A > 0 and yU > 0. Then, zU = 39.4 is a possible
value of z, and we can go further in the computation, as follows
1.8 An Example of Solution of Lambert’s Problem Using Universal Variables 67
12 12
yU 3:9211363 106
xU ¼ ¼ ¼ 2:8145 106
CU 4:95 107
At x = xU and within the interval zL < z < zU, we seek a zero of the following
function
1
1
f ðDt0 ; xU Þ ¼ 3600 x3U SU þ Ay2U 1 ¼ 3600 ð22:295 1018 0:031
l2
1
þ 1:23723 104 1:9802 103 Þ 1 1015
6:31348 102
The value of f(Dt0, xU) is strongly negative. Consequently, if z = zU, we are

largely away from the point x in which the function f(Dt0, x) vanishes. The sign in
front of f(Dt0, xU) and the sign in front of f(Dt0, xL) are important, because, in case of
f(Dt0, x) having opposite signs at the endpoints xL and xU, there is assurance of the
existence of a zero of f(Dt0, x) between xL and xU.
Now, to determine the sign of y at the lower endpoint (zL) of the interval, we
compute the values of the Stumpff functions C(z) and S(z) at z = zL, as follows
1

1 cosh 39:42
CL ¼ C ð39:4Þ ¼ ¼ 6:7279
ð39:4Þ
1
1
sinh 39:42 39:42
SL ¼ Sð39:4Þ ¼ 3 ¼ 1:0505
39:42

1
yL ¼ ri þ rf þ AðzL SL 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103
CL2
1
þ 1:23723 104 ð39:4 1:0505 1Þ 1 ¼ 17:4436 10
4
ð6:7279Þ2
Thus, at z = zL, there results A > 0 and yL < 0. Then, zL = −39.4 is not a pos-
sible value of z. Consequently, zL must be increased to a value greater than −39.4,
in order that the new value of zL should lead to a positive value of yL.
Trying with values of zL which are progressively higher than −39.4 and checking
the sign of the correspondent values of yL, we find
zL ¼ 3:1
which satisfies the condition A > 0 and yL > 0, as will be shown below. Thus, we
discard the initial value zL = −39.4 and compute the functions C(z) and S(z) at the
new lower endpoint (zL = −3.1) of the interval zL < z < zU, as follows
1
1 cosh 3:12
CL ¼ C ð3:1Þ ¼ ¼ 0:64328
ð3:1Þ
1 1
sinh 3:12 3:12
SL ¼ Sð3:1Þ ¼ 3 ¼ 0:194492
3:12

1
yL ¼ ri þ rf þ AðzL SL 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103
CL2
1
þ 1:23723 104 ð3:1 0:194492 1Þ 1
ð0:64328Þ2
¼ 0:303254 104
Now, there results A > 0 and yL > 0; consequently, we can compute

12 12
yL 0:303254 104
xL ¼ ¼ ¼ 0:686599 102
CL 0:64328
1
1
f ðDt0 ; xL Þ ¼ 3600 x3L SL þ Ay2L 1 ¼ 3600 ð0:323675 106 0:194492
l2
1
þ 1:23723 10 0:550685 102 Þ
4
¼ 2:42113 103
6:31348 102
Since there results f(Dt0, xL) > 0, then f(Dt0, x) = 0 has a real root for x falling
between xL and xU, that is,
0:686599 102 x 2:8145 106
or, which is the same, f(Dt0, z) = 0 has a real root for z falling between zL and zU,
that is, −3.1 z 39.4.
On the other hand, as has been shown above, the value zU = 39.4 leads to a large
negative value of f(Dt0, xU). Consequently, the root z cannot be in the vicinity of
39.4. We try zU = 5 (i.e., −3.1 z 5) and check whether f(Dt0, x) has opposite
signs at these two endpoints. Since zL = −3.1 leads to f(Dt0, xL) = 2.42113 103,
then we set zU = 5 and compute the Stumpff functions
1
1 cos 52
CU ¼ C ð5Þ ¼ ¼ 0:32345
5
1 1
52 sin 52
SU ¼ Sð 5Þ ¼ 3 ¼ 0:12963
52

1
yU ¼ ri þ rf þ AðzU SU 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103
CU2
1
þ 1:23723 104 ð5 0:12963 1Þ 1 ¼ 2:01048 10
4
ð0:32345Þ2
There results A > 0 and yU > 0; then, we can compute

12 12
yU 2:01048 104
xU ¼ ¼ ¼ 2:493139 102
CU 0:32345
The function f(Dt0, xU) is then

1
1
f ðDt0 ; xU Þ ¼ 3600 x3U SU þ Ay2U 1 ¼ 3600 ð15:49671 106 0:12963
l2
1
þ 1:23723 104 1:417915 102 Þ ¼ 2:36 103
6:31348 102
Since the signs of f(Dt0, xL) and f(Dt0, xU) are opposite, then the unknown root
z of f(Dt0, z) = 0 is between −3.1 and 5. Unlike the authors cited previously [18, 19,
and 21], we use Müller’s method of parabolic interpolation, which has been
described in Sect. 1.5, to find iteratively the unknown value of z. This method
operates a quadratic interpolation and consequently requires three points in the
vicinity of the root to be found.
First iteration. We set
z2 ¼ 3:1
z1 ¼ 5
so that z2 should be less than z1. The respective values of the time function f,
computed above, are
f2 ¼ f ð3600; 3:1Þ ¼ 2:42113 103

f1 ¼ f ð3600; 5Þ ¼ 2:36 103
Now, we choose the third point so that its abscissa z0 should be placed between
z2 and z1, that is, so that
z2 \z0 \z1
A possible way to do this is to bisect the interval, by taking z0 = (z2 + z1)/2; for
convenience, we choose arbitrarily z0 = 1, which value satisfies the condition
3:1\z0 \5
Then we compute the time function f0 = f(3600, 1) corresponding to z0 = 1, as

shown below:
1
1 cos 12
C0 ¼ Cðz0 Þ ¼ Cð1Þ ¼ ¼ 0:4596977
1
1 1
12 sin 12
S0 ¼ Sð x 0 Þ ¼ Sð 1Þ ¼ 3 ¼ 0:158529
12

1
y0 ¼ ri þ rf þ Aðz0 S0 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103
C02
1
þ 1:23723 104 ð1 0:158529 1Þ 1
ð0:4596977Þ2
¼ 2:77591 104
There results A > 0 and y0 > 0; therefore, we can compute

12 1
y0 2:77591 104 2
x0 ¼ ¼ ¼ 1:6426468 102
C0 0:4595977
The function f0 f(Dt0, x0) is then

1
1
f ðDt0 ; x0 Þ ¼ 3600 x30 S0 þ Ay20 1 ¼ 3600 ð4:432335 106 0:158529
l2
1
þ 1:23723 104 1:1137311 102 Þ
6:31348 102
¼ 0:30452 103
Starting from the values

f2 ¼ 2:42113 103 s
f0 ¼ 0:30452 103 s
f1 ¼ 2:36 103 s
obtained above, we compute
h1 ¼ z1 z0 ¼ 5 1 ¼ 4
h2 ¼ z0 z2 ¼ 1 þ 3:1 ¼ 4:1
c ¼ h2 =h1 ¼ 4:1=4 ¼ 1:025
and then the three coefficients (a, b, and c) of the quadratic polynomial which
interpolates the three points (z2, f2), (z0, f0), and (z1, f1), as follows
a¼ ¼ 0:018504
ch21 ð1 þ cÞ
f1 f0 Ah21
b¼ ¼ 666:055984
h1
c ¼ f0 ¼ 0:30452 103
Thus, z is estimated as follows
2c
z ¼ z0 1 ¼ 1:45719
b ð b2 4acÞ2
where the sign in front of the square root of the discriminant b2 − 4ac has been set
to minus because b has a negative value, as has been shown in Sect. 1.5.
Now, we compute the Stumpff functions C(z) and S(z) at z = 1.45719, as follows
1

1 cos 1:457192
C ðzÞ ¼ Cð1:45719Þ ¼ ¼ 0:4421547
1:45719
1 1
1:457192 sin 1:457192
SðzÞ ¼ Sð1:45719Þ ¼ 3 ¼ 0:1549361
1:457192

1
y ¼ ri þ rf þ AðzS 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103 þ 1:23723
C2
1
10 ð1:45719 0:1549361 1Þ
4
1
ð0:4421574Þ2
¼ 1:33536 104
There results A > 0 and y > 0; then, we can compute
y 12 12
1:33536 104
x¼ ¼ ¼ 1:73784 102
C 0:4421574
The function f f(Dt0, x) is then

1
1
f ðDt0 ; xÞ ¼ 3600 x3 S þ Ay2 1 ¼ 3600 ð5:24846 106 0:15494
l2
1
þ 1:23723 104 1:15558 102 Þ
6:31348 102
¼ 47:42 s
Thus, taking z = 1.45719 leads to f(Dt0, x) = 47.42 instead of f(Dt0, x) = 0.

Since z = 1.45719 is greater than z0 = 1, then z0, z2, and z are taken for the next
iteration.
Second iteration.
We reset the subscripts 0, 1, and 2 so that z2 < z0 < z1. In other words, we set
z2 ¼ 3:1
z0 ¼ 1:0
z1 ¼ 1:45719
It follows that
f2 ¼ f ðz2 Þ ¼ 2:42113 103 s

f0 ¼ f ðz0 Þ ¼ 0:30452 103 s
f1 ¼ f ðz1 Þ ¼ 47:42 s
Then, we compute again
h1 ¼ z1 z0 ¼ 1:45719 1:0 ¼ 0:45719

h2 ¼ z0 z2 ¼ 1:0 þ 3:1 ¼ 4:1
c ¼ h2 =h1 ¼ 4:1=0:45719 ¼ 8:96783
interpolates the three points (z2, f2), (z0, f0), and (z1, f1), as follows
a¼ ¼ 10:116502
ch21 ð1 þ cÞ
f1 f0 Ah21
b¼ ¼ 557:72309
h1
c ¼ f0 ¼ 0:30452 103
2c
z ¼ z0 1 ¼ 1:5407026
b ðb2 4acÞ2
Now, we compute the Stumpff functions C(z) and C(z) at z = 1.5407026, as

follows
1

1 cos 1:54070262
C ðzÞ ¼ Cð1:5407026Þ ¼ ¼ 0:43901
1:5407026
1 1
1:54070262 sin 1:54070262
SðzÞ ¼ Sð1:5407026Þ ¼ 3 ¼ 0:15429
1:54070262

1
y ¼ ri þ rf þ AðzS 1Þ 1 ¼ ½11:37585 103 þ 16:38322 103 þ 1:23723
C2
1
10 ð1:5407026 0:15429 1Þ
4
1 ¼ 1:3525 10
4
ð0:43901Þ 2
There results A > 0 and y > 0; thus, we can compute
y 12 12
1:3525 104
x¼ ¼ ¼ 1:75522 102
C 0:43901
The function f f(Dt0, x) is then

1
1
f ðDt0 ; xÞ ¼ 3600 x3 S þ Ay2 1 ¼ 3600 ð5:40746 106 0:15429
l2
1
þ 1:23723 104 1:16297 102 Þ ¼ 0:5 s
6:31348 102
Thus, taking z = 1.5407026 leads to f(Dt0, x) = −0.5 instead of f(Dt0, x) = 0.

Since the relative error amounts to half a second out of 3600 s, we proceed no
further with the iterations and take z = 1.5407026 as the correct value. The cor-
respondent value of y, as shown above, is 1.3525 104.
Thus, the Lagrangian coefficients f, g, and g′ are
y 1:3525 104
f ¼1 ¼1 ¼ 0:18892
ri 1:137585 104
y 1:3525 104
g0 ¼ 1 ¼ 1 ¼ 0:17446
rf 1:638322 104
12 1
y 1:3525 104 2
g¼A ¼ 1:23723 104 ¼ 0:22790 104
l 3:986 105
The three components of the velocity vector at the initial point of the transfer
trajectory, with respect to the geocentric-equatorial reference system, are
Xf fXi 14:6 103 þ 0:18892 5 103

Xi0 ¼ ¼ ¼ 5:992 km/s
g 0:22790 104
Yf fYi 2:5 103 þ 0:18892 10 103
Yi0 ¼ ¼ ¼ 1:926 km/s
g 0:22790 104
Zf fZi 7 103 þ 0:18892 2:1 103
Zi0 ¼ ¼ ¼ 3:246 km/s
g 0:22790 104
The three components of the velocity vector at the final point of the transfer
trajectory, with respect to the same reference system, are
g0 Xf Xi 0:17446 14:6 103 5 103

Xf0 ¼ ¼ ¼ 3:312 km/s
g 0:22790 104
g0 Yf Yi 0:17446 2:5 103 10 103
Yf0 ¼ ¼ ¼ 4:197 km/s
g 0:22790 104
g0 Zf Zi 0:17446 7 103 2:1 103
Zf0 ¼ ¼ ¼ 0:386 km/s
g 0:22790 104
These results may be checked against those obtained by Curtis [21], who uses
MATLAB and chooses a tolerance of 1 10−8:
Xi0 ¼ 5:99249 km=s Xf0 ¼ 3:31246 km=s

Yi0 ¼ 1:92536 km=s Yf0 ¼ 4:19662 km=s
0
Zi ¼ 3:24564 km=s Zf0 ¼ 0:385288 km=s
The readers who have a computer and a FORTRAN compiler can solve the
problem of determining iteratively y (and hence the initial and final velocity vec-
tors) by using a subroutine, available in [22], based on Müller’s method. Directions
on the use of this subroutine are given in [23].
Vallado [17] has written a MATLAB function (Lambertu.m), for the solution of
Lambert’s problem in terms of universal variables. This function is also available in
[24].
The transfer trajectory computed above turns out to be an ellipse, as will be
shown below. First, the magnitude of the velocity vector at the initial point is
1 1
v ri0 ¼ Xi02 þ Yi02 þ Zi02 2 ¼ 5:9922 þ 1:9262 þ 3:2462 2 ¼ 7:082 km/s
On the other hand, there results

12 1
2l 2 3:986 105 2
¼ ¼ 8:371 km/s
ri 11:376 103
Since the values of vi and (2l/ri)½ computed above are such that
1
2l 2
vi \
ri
then the conic transfer trajectory is an ellipse, as has been shown in Sect. 1.5.
1.9 The Classical Orbital Elements
As has been shown in Sect. 1.1, the trajectory followed by a particle, which is acted
upon only by the gravitational force due to another particle, is a conic section. The
parameters which must be specified to identify uniquely this conic section and the
position of the particle moving along it are called the orbital elements, which are
shown in the following figure (courtesy of NASA [25]).
In order to describe the motion of a particle orbiting in the gravitational field

exerted by another particle, it is necessary to specify six quantities: the three
components of the position vector and the three components of the velocity vector
of the orbiting particle with respect to a reference system having its origin in the
centre of mass of the two particles.
Five quantities are necessary to specify the size, shape, and orientation of an
orbit. A sixth quantity is necessary to specify the position of the particle along its
orbit at any given time. Two (a and e) of these quantities define the shape of the
orbit; three (X, i, and x) define the orientation of the orbit in space; and one (/0)
defines the position of the orbiting particle at a given time t0.
The preceding figure refers to an artificial satellite revolving about the Earth. The
geocentric-equatorial reference system XYZ has its origin in the centre of the Earth,
and its X-axis points from the Earth towards the first point of Aries (indicated by the
symbol ♈), that is, towards the position of the Sun at the vernal equinox. This
direction is the intersection of the equatorial plane of the Earth and the ecliptic
plane, which is the plane containing the orbit of the Earth around the Sun.
Therefore, the X-axis lies in both of these planes. The Z-axis is parallel to the
rotation axis of the Earth and points towards the north pole. The Y-axis completes
the right-handed orthogonal set XYZ.
1.9 The Classical Orbital Elements 77
In other words, let uX, uY, and uZ be the unit vectors relating to the axes,
respectively, X, Y, and Z. The following equality holds:
uY ¼ uZ uX
The XYZ reference system does not rotate with the Earth; on the contrary, the
Earth rotates with respect to this system, which is, to the first order, fixed with
respect to the distant (also called fixed) stars. This is the reference system com-
monly used in satellite orbit calculations. It is to be noted that XYZ is subject to
second-order changes with time, owing to various slow motions of the rotation axis
of the Earth with respect to the fixed stars. In order to take account of this fact, the
geocentric-equatorial co-ordinates XYZ are associated with a date (also called
epoch) to which these co-ordinates apply. Spacecraft orbits and positions are often
expressed for the standard astronomical epoch known as J2000.0 (the 1st of January
2000 at noon at the Greenwich meridian).
The six classical (also called Keplerian) orbital elements are:
• major semi-axis (a);
• eccentricity (e);
• inclination (i) of the orbital plane with respect to the equator, or, which is the
same, the angle between the unit vector uZ and the angular momentum per unit
mass (h);
• right ascension of the ascending node (X), which is the angle, measured
counterclockwise in the fundamental plane by an observer who is in the northern
semi-space, between the unit vector uX and the direction (line of nodes) of the
point where the satellite crosses the fundamental plane from south to north;
• argument of perigee (x), which is the angle, measured in the direction of motion
and in the orbital plane, between the ascending node and the perigee; and
• time of perigee passage (tp), or the time at which the satellite is at perigee.
Another orbital element used instead of the major semi-axis (a) is the semi-latus
rectum (p), which depends on a and e through p = a(1 − e2).
Another orbital element used instead of the argument of perigee (x) is the
longitude of perigee (P), which is the angle from the unit vector uX to the perigee
direction, measured counterclockwise in the orbital plane by an observer who is in
the semi-space of positive h.
Other orbital elements commonly used instead of the time of perigee passage
are:
• true anomaly at epoch (/0), which is the angle, measured in the orbital plane,
between perigee and the position of the satellite at a given time t0, called the
epoch;
• mean anomaly at epoch (M0), which is the value of the mean anomaly M,
defined in Sect. 1.3, at epoch (t0);
• argument of latitude at epoch (u0), which is the angle, measured in the orbital
plane, between the ascending node and the position vector r0 of the satellite at
epoch;
• true longitude at epoch (‘0), which is the angle between the unit vector uX and
the position vector r0 of the satellite at epoch.
The following figure (courtesy of NASA [25]) shows some of these angles.
It is to be noted that some of the elements indicated above may be undefined in

particular cases. For example, in case of a circular orbit, the perigee is undefined,
and therefore x is also undefined. In case of an equatorial orbit, the line of nodes is
undefined, and therefore X is also undefined. If both X and x are defined, that is, in
case of an elliptic non-equatorial orbit, then the definition of P given above implies
that P = X + x; otherwise P is undefined.
If both x and /0 are defined, then u0 = x + /0. In case of an equatorial orbit,
both x and u0 are undefined. If X, x, and /0 are all defined, then
‘ 0 ¼ X þ x þ / 0 ¼ P þ / 0 ¼ X þ u0
In case of an equatorial orbit, then ‘0 = P + /0. In case of a circular orbit, then

there results ‘0 = X + u0. In case of an orbit being both circular and equatorial, then
‘0 is the angle from the unit vector uX to the position vector (r0) of the satellite at
epoch. Since uX and r0 are always definite, then ‘0 is always definite.
Given the position vector (r) and velocity vector (r′) of a satellite with respect to
the geocentric-equatorial reference system XYZ at a particular time t, it is possible to
compute the six orbital elements, as the following example will show. To this end,
we use the data of the transfer ellipse, at time t = ti, given in Sect. 1.8. They are
Xi ¼ 5 103 km Xi0 ¼ 5:992 km=s

Yi ¼ 10 103 km Yi0 ¼ 1:926 km=s
Zi ¼ 2:1 103 km Zi0 ¼ 3:246 km=s
Let
1 1
ri ¼ Xi2 þ Yi2 þ Zi2 2 ¼ 50002 þ 100002 þ 21002 2 ¼ 11376 km
1 1
v ri0 ¼ Xi02 þ Yi02 þ Zi02 2 ¼ 5:9922 þ 1:9262 þ 3:2462 2 ¼ 7:082 km/s
be, respectively, the magnitude of the position vector (ri) and the magnitude of the
velocity vector (r′i) of the satellite at t = ti. The major semi-axis of the transfer
ellipse can be computed by means of the vis-viva integral
v2i 2 1
¼
l ri a
This equation, solved for the major semi-axis, yields

1 1
a¼ ¼ ¼ 2:0 104 km
2 v2i 2 7:0822

ri l 11376 3:986 105
The moment-of-momentum vector per unit mass (h = ri r′i), which is a

constant of motion in the two-body problem, is computed as follows

h ¼ hX uX þ hY uY þ hZ uZ ¼ Yi Zi0 Zi Yi0 uX þ Zi Xi0 Xi Zi0 uY

þ Xi Yi0 Yi Xi0 uZ ¼ 10 103 3:246 2:1 103 1:926 uX

þ 2:1 103 5:992 5 103 3:246 uY

þ 5 103 1:926 þ 10 103 5:992 uZ
¼ 28:4154 103 uX 28:8132 103 uY þ 69:55 103 uZ
The magnitude of h is
1 1
h ¼ ðh hÞ2 ¼ 28415:42 þ 28813:22 þ 695502 2 ¼ 80466:478 km/s
The semi-latus rectum results from
h2 80466:4782
p¼ ¼ ¼ 16243:949 km
l 398600
The scalar eccentricity results from


p ¼ a 1 e2
which, solved for e, yields
1
p12 16243:949 2
e¼ 1 ¼ 1 ¼ 0:43336
a 20000
The following equation derived in Sect. 1.1
r0i h ri
e¼
l ri
makes it possible to compute the eccentricity vector (e). Since ri and r′i are given
and h = ri r′i has been computed above, then their respective values can be
substituted in the preceding equation. By so doing, we find
r0i h ri
e¼ ¼ 0:131178 uX þ 0:3978756 uY þ 0:1112378 uZ
l ri
As a check, the scalar eccentricity e may be computed again by evaluating the

magnitude of the eccentricity vector. This yields
1 1
e ¼ ðe eÞ2 ¼ 0:1311782 þ 0:39787562 þ 0:11123782 2 ¼ 0:43346
The orbital plane is inclined with respect to the equatorial plane by

hZ 69:55 103
i ¼ arccos ¼ arccos ¼ 30 :194
h 80:466378 103
The right ascension of the ascending node (X) results from computing first the
node vector (n) defined by
n uZ h
Thus, there results

n ¼ nX uX þ nY uY ¼ hY uX þ hX uY ¼ 2:1 103 5:992 5 103 3:246 uX

þ 10 103 3:246 2:1 103 1:926 uY

¼ 28:8132 103 uX þ 28:4154 103 uY
The magnitude of the node vector is

1 1
n ¼ n2X þ n2Y 2 ¼ 28813:22 þ 28;415:42 2 ¼ 40467:709
This done, the angle X results from

n
X
X ¼ arccos ð nY 0Þ
n n
X ¼ 360 arccos
X
ðnY \0Þ
n
In the present case, nY = 28.4154 103 > 0; hence

n
28813:2
¼ 44 :602
X
X ¼ arccos ¼ arccos
n 40467:709
The argument of perigee (x) results from

n e
X ¼ arccos ðeZ 0Þ
ne
n e
X ¼ 360 arccos ðeZ \0Þ
ne
In the present case, eZ = 0.1112378 > 0; hence

n e
28813 0:13118 þ 28415 0:39788
X ¼ arccos ¼ arccos ¼ 30 :683
ne 40468 0:43346
The true anomaly at epoch (/0) results from

e ri
/0 ¼ arccos ri r0i 0
eri

e ri
/0 ¼ 360 arccos ðri r0i \0Þ
eri
ri r0i ¼ 5 103 ð5:992Þ þ 10 103

1:926 þ 2:1 103 3:246 ¼ 3883:4\0
Therefore,

0:13118 5000 þ 0:39788 10000 þ 0:11124 2100
/0 ¼ 360 arccos
0:43346 11376
¼ 350 :85
Montenbruck and Gill [26] suggest another way to compute the argument of
perigee and the true anomaly at epoch. This way will be shown below with the view
of checking the values computed previously.
The argument of latitude at epoch (u0) results from
2 3
6 Z 7
u0 ¼ arctan6
4 i 7
hY hX 5
Xi þ Yi
h h
Montenbruck and Gill point out that in expressions of the form

y
a ¼ arctan
x
the quadrant of a must be chosen so that the sign of the denominator (x) should be
equal to the sign of cos a. Since cos a is positive in the first and fourth quadrants
and negative in the second and third, then a is between −90° and +90° for positive
values of x, and between +90° and +270° for negative values of x. In the present
case, there results
hX Yi Zi0 Zi Yi0 10000 3:246 2100 1:926

¼ ¼ ¼ 0:35313
h h 80466
0 0
hY Zi Xi Xi Zi 2100 5:992 5000 3:246
¼ ¼ ¼ 0:35808
h h 80466
Therefore, the argument of latitude at epoch is

2100
u0 ¼ arctan ¼ 21 :53
5000 ð0:35808Þ þ 10000 0:35313
The mean motion (n) results from
l 12 1
398600 2
n¼ ¼ ¼ 2:23215 104 s1
a3 200003
The eccentric anomaly at epoch (Æ0) results from

0 1 0 1
ri r0i 3883:4
B 2 C B200002 2:23215 104 C
Æ 0 ¼ arctan@ a nri A ¼ arctanB
@
C ¼ 354 :24
A
1 11376
a 1
20;000
The true anomaly at epoch (/0) results from

" 1 #
ð1 e2 Þ2 sin Æ 0
/0 ¼ arctan
cos Æ 0 e
" 1
#
ð1 0:433462 Þ2 sin 354 :24
¼ arctan ¼ 350 :85
cos 354 :24 0:43346
The argument of perigee (x) results from
x ¼ u0 /0 ¼ 21 :53 ð9 :15Þ ¼ 30 :68
For convenience of the reader, the orbital elements computed above are shown
on the left-hand side of the following table. They can be compared with those
computed by Curtis [21] by means of the MATLAB function Lambert (Example
5.2), which are shown on the right-hand side of the same table.
a ¼ 2:0 104 km a ¼ 2:00029 104 km

e ¼ 0:43336 e ¼ 0:433488
i ¼ 30 :193 i ¼ 30 :191
X ¼ 44 :602 X ¼ 44 :6002
x ¼ 30 :683 x ¼ 30 :7062
/0 ¼ 350 :85 /0 ¼ 350 :83
Let us consider now the inverse problem. Given a set of six independent orbital
elements, we want to compute the three components of the position vector r and the
three components of the velocity vector r′ at some epoch t0. To this end, we proceed
in two steps: first, r and r′ are expressed in the perifocal reference system Oxyz,
which is shown in the following figure.
Then, r and r′ are expressed in the geocentric-equatorial reference system OXYZ.

The perifocal reference system has its origin O in the centre of mass of the principal
attracting body. For example, in case of an artificial satellite revolving about the
Earth, the origin of the perifocal reference system is the centre of the Earth; in case
of a planet revolving about the Sun, the origin of this system is the centre of the
Sun.
The fundamental plane of the perifocal reference system is the plane to which the
motion of the orbiting body is confined. Let ux, uy, and uz be the three unit vectors
along, respectively, the x, y, and z axes of the perifocal reference system. Of these
unit vectors, ux and uy are along the two Cartesian axes x and y which are contained
in the orbital plane, oriented so that ux points towards the perifocus (or perigee, for
an artificial satellite of the Earth) and uy is 90° ahead of ux in the direction of orbital
motion, whereas uz is the unit vector along z, perpendicular to the orbital plane,
oriented so as to make xyz a right-handed system of reference. To this end, uz points
towards the same direction as the moment of momentum per unit mass h = r r′
of the satellite.
Given the values of the six classical elements a, e, i, X, x, and /, the position
vector r, with respect to the perifocal reference system, is expressible as follows
r ¼ ðr cos /Þux þ ðr sin /Þuy
where the magnitude (r) of r results from the equation of a conic section
að 1 e 2 Þ
r¼
1 þ e cos /
In order to express the velocity vector r′ with respect to the same (perifocal)
reference system, we differentiate the equation r = (r cos /)ux + (r sin /)uy with
respect to time, bearing in mind that ux and uy are, both of them, constant with time,
so that u′x = u′y = 0. This yields
r0 ¼ ðr 0 cos / r/0 sin /Þux þ ðr 0 sin / þ r/0 cos /Þuy
Now, remembering that h = r2/′, that p = a(1 − e2) = h2/l, and differentiating
the equation of a conic section written above with respect to time, there results
2 1
!
0 p 0 r 0 h ð l pÞ 2
r ¼ 2
ðe/ sin /Þ ¼ e/ sin / ¼ e sin / ¼ e sin /
ð1 þ e cos /Þ p p p
12
l
¼ e sin /
p
and
12
h 1 þ e cos / l
r/0 ¼ ¼h ¼ ð1 þ e cos /Þ
r p p
By substituting the two expressions written above into
r0 ¼ ðr 0 cos / r/0 sin /Þux þ ðr 0 sin / þ r/0 cos /Þuy
and simplifying, there results

12
l

r0 ¼ ðsin /Þux þ ðe þ cos /Þuy
p
In summary, in the perifocal reference system Oxyz, the position vector r and the
velocity vector r′ can be expressed as a function of the classical orbital elements as
follows

12
0 l

r ¼ ðsin /Þux þ ðe þ cos /Þuy
p
where

p ¼ a 1 e2
p
r¼
1 þ e cos /
It remains to show how r and r′ are to be transformed in order express them, as a

function of the classical orbital elements, in the geocentric-equatorial reference
system OXYZ.
This transformation can be done by means of the angles X, x, and i, as will be
shown below. Let ux, uy, and uz be the unit vectors of the perifocal reference system
Oxyz, having its origin (O) in the centre of the Earth. Let uX, uY, and uZ be the unit
vectors of the geocentric-equatorial reference system OXYZ, having the same origin
O. The angles X, x, and i define the orientation of the perifocal reference system
Oxyz with respect to the geocentric-equatorial reference system OXYZ. These angles
are shown in the following figure, which is due to the courtesy of Wikimedia [27].
Let rij be the element of the ith row and the jth column of the following matrix
2 3
uX ux uX uy uX uz
R 4 uY ux uY uy uY uz 5
uZ ux uZ uy uZ uz
The nine elements of this matrix are the direction cosines of the axes x, y, and
z of the perifocal system with respect to the geocentric-equatorial system XYZ.
These elements can be determined as follows. Let us consider two (R1 and R3) of
the following three elementary rotation matrices
2 3 2 3
1 0 0 cos a 0 sin a
6 7 6 7
R1 4 0 cos a sin a 5 R2 4 0 1 0 5
0 sin a cos a sin a 0 cos a
2 3
cos a sin a 0
6 7
R3 4 sin a cos a 0 5
0 0 1
With reference to the preceding figure, the axes x, y, and z of the perifocal system
are brought to coincide with the axes, respectively, X, Y, and Z of the geocentric-
equatorial system by means of the following sequence of elementary rotations:
• first, a rotation through the angle −x about the z-axis;
• then, a rotation through the angle −i about the new x-axis, which is the line of
nodes;
• finally, a rotation through the angle −X about the new z-axis, which is the Z-
axis.
In mathematical terms, the matrix R results from the following product

R ¼ R3 ðXÞR1 ðiÞR3 ðxÞ
where the minus sign in front of the angles x, i, and X is necessary because the
three rotations occur clockwise. The resulting matrix R is

cos X cos x sin X sin x cos i cos X sin x sin X cos x cos i sin X sin i

sin X cos x þ cos X sin x cos i sin X sin x þ cos X cos x cos i cos X sin i

sin x sin i cos x sin i cos i
Having determined the components of the rotation matrix R, the components X,

Y, and Z of the position vector r in the geocentric-equatorial reference system and
the components X′, Y′, and Z′ of the velocity vector r′ in the same reference system
result from
2 3 2 3 2 03 2 03
X x X x
4 Y 5¼ R4 y 5 4 Y 0 5¼ R4 y0 5
Z 0 Z0 0
where
12
l
x ¼ r cos / x0 ¼ sin /
p
12
l
y ¼ r sin / y0 ¼ ðe þ cos /Þ
p
are the components of, respectively, the position vector and the velocity vector in
the perifocal reference system, and
p
p ¼ a 1 e2 r¼
1 þ e cos /
are, respectively, the semi-latus rectum and the radius vector at epoch.
As is easy to verify, R is an orthogonal matrix, because it satisfies the condition
RT R ¼ I
where RT is the transpose of R, and I is the following 3 3 identity matrix

2 3
1 0 0
I 40 1 05
0 0 1
Since the inverse matrix R−1 must, by definition, satisfy the condition
R1 R ¼ I
then the orthogonality of R implies that R−1 = RT.

Consequently, in order to transform geocentric-equatorial co-ordinates XYZ into
perifocal co-ordinates xyz, we can use the transpose RT, whose components are

cos X cos x sin X sin x cos i sin X cos x þ cos X sin x cos i sin X sin i

cos X sin x sin X cos x cos i sin X sin x þ cos X cos x cos i cos X sin i

As an example of application of the formulae indicated above, given the orbital

elements of an artificial satellite revolving about the Earth, which are
a ¼ 2:0 104 km X ¼ 44 :6017

e ¼ 0:43336 x ¼ 30 :68
i ¼ 30 :193 /0 ¼ 350 :85
we want to compute the three components X, Y, and Z of the position vector r and
the three components X′, Y′, and Z′ of the velocity vector r′ in the
geocentric-equatorial reference system, at some epoch t0.
As has been shown above, we compute first the semi-latus rectum p and the
radius vector r0 of the satellite at t = t0, as follows.

p ¼ a 1 e2 ¼ 2:0 104 1 0:433362 ¼ 1:624398 104 km
p 1:624398 104
r0 ¼ ¼ ¼ 1:1376 104 km
1 þ e cos /0 1 þ 0:43336 cos 350 :85
This done, we compute the two components x and y of the position vector r and
the two components x′ and y′ of the velocity vector r′ in the perifocal reference
system Oxyz having its origin O in the centre of the Earth, as follows
x ¼ r0 cos /0 ¼ 1:1376 104 cos 350 :85 ¼ 1:1232 104 km

y ¼ r0 sin /0 ¼ 1:1376 104 sin 350 :85 ¼ 0:1809 104 km
1
1
x0 ¼ ðl=pÞ =2 sin /0 ¼ 3:986 105 = 1:624398 104 2 sin 350 :85 ¼ 0:7878 km=s
1
1
y0 ¼ ðl=pÞ =2 ðe þ cos /0 Þ ¼ 3:986 105 = 1:624398 104 2 ð0:4334588
þ cos 350 :85Þ ¼ 7:0382 km=s
Now, we use the components of the rotation matrix R determined above to

compute the three components X, Y, and Z of the position vector r and the three
components X′, Y′, and Z′ of the velocity vector r′ in the geocentric-equatorial

reference system OXYZ, as follows
X ¼ r11 x þ r12 y ¼ ðcos X cos x sin X sin x cos iÞx þ ð cos X sin x
sin X cos x cos iÞy ¼ ðcos 44 :6017 cos 30 :68 sin 44 :6017
sin 30 :68 cos 30 :193Þ 1:1232 104 þ ð cos 44 :6017 sin 30 :68

sin 44 :6017 cos 30 :68 cos 30 :193Þ 0:1809 104 ¼ 0:5 104 km
Y ¼ r21 x þ r22 y ¼ ðsin X cos x þ cos X sin x cos iÞx þ ð sin X sin x
þ cos X cos x cos iÞy ¼ ðsin 44 :6017 cos 30 :68 þ cos 44 :6017
sin 30 :68 cos 30 :193Þ 1:1232 104 þ ð sin 44 :6017 sin 30 :68

þ cos 44 :6017 cos 30 :68 cos 30 :193Þ 0:1809 104 ¼ 1:0 104 km
Z ¼ r31 x þ r32 y ¼ ðsin x sin iÞx þ ðcos x sin iÞy ¼ ðsin 30 :68 sin 30 :193Þ

1:1232 104 þ ðcos 30 :68 sin 30 :193Þ 0:1809 104 ¼ 0:21 104 km
X 0 ¼ r11 x0 þ r12 y0 ¼ ðcos X cos x sin X sin x cos iÞx0 þ ð cos X sin x
sin X cos x cos iÞy0 ¼ ðcos 44 :6017 cos 30 :68 sin 44 :6017
sin 30 :68 cos 30 :193Þ 0:7878 þ ð cos 44 :6017 sin 30 :68
sin 44 :6017 cos 30 :68 cos 30 :193Þ 7:0382 ¼ 5:992 km=s
Y 0 ¼ r21 x0 þ r22 y0 ¼ ðsin X cos x þ cos X sin x cos iÞx0 þ ð sin X sin x
þ cos X cos x cos iÞy0 ¼ ðsin 44 :6017 cos 30 :68 þ cos 44 :6017
sin 30 :68 cos 30 :193Þ 0:7878 þ ð sin 44 :6017 sin 30 :68
þ cos 44 :6017 cos 30 :68 cos 30 :193Þ 7:0382 ¼ 1:926 km=s
Z 0 ¼ r31 x0 þ r32 y0 ¼ ðsin x sin iÞx0 þ ðcos x sin iÞy0 ¼ ðsin 30 :68 sin 30 :193Þ 0:7878
þ ðcos 30 :68 sin 30 :193Þ 7:0382 ¼ 3:246 km=s
These values turn out to be the same as those which were given as input data in
the preceding example.
1.10 Orbital Elements Defined for Any Orbit
As has been shown in Sect. 1.9, some of the classical orbital elements are undefined
in certain cases. The line of nodes is undefined for orbits of zero inclination angle
(i = 0°). The line of apsides is undefined for circular orbits (e = 0). The orbital
elements affected in these cases are the right ascension of the ascending node (X),
the argument of perigee (x) and the time of perigee passage (tp). To avoid these
difficulties, some authors have proposed to use alternative sets of orbital elements
which are free from problems arising with orbits of small eccentricity or inclination.
One of these sets uses the so-called equinoctial elements, which are defined below
a q1 ¼ tanði=2Þ sin X
p1 ¼ e sinðX þ xÞ q2 ¼ tanði=2Þ cos X
p2 ¼ e cosðX þ xÞ l0 ¼ M0 þ X þ x
where a is the major semi-axis, M0 is the mean anomaly, and l0 is called the mean
longitude at epoch. The name “equinoctial elements” is due to Arsenault et al. [28].
In order to avoid confusion in the nomenclature, we use here, in accordance with
Battin [11], the symbols p1 and p2 for, respectively, the second and the third element,
and the symbols q1 and q2 for, respectively, the fourth and the fifth element.
Comparing the equinoctial elements to the classical elements, the first element of
the former set (a) has the same meaning as that of the latter. The second and the
third equinoctial element (p1 and p2) are the two components of the eccentricity
vector (e), which lies in the orbital plane and points from the centre of the Earth to
perigee. The fourth and the fifth equinoctial element (q1 and q2) are the two
components of the ascending node vector (n), which points from the centre of the
Earth to the ascending node.
Consequently, the scalar eccentricity (e) is expressible as a function of p1 and p2,
as follows
1
e ¼ p21 þ p22 2
The orbital inclination (i) is expressible as a function of q1 and q2, as follows

h 1 i
i ¼ 2 arctan q21 þ q22 2
When the orbital inclination is greater than a given minimum value (e.g., greater
than e = 1 10−6), then the right ascension of the ascending node (X) is
expressible as follows

q1
X ¼ arctan
q2
Otherwise, in case of an equatorial orbit, X is undetermined.

Likewise, when the orbital eccentricity is greater than a given minimum value,
then the argument of perigee (x) is expressible as follows

p1
x ¼ arctan X
p2
Otherwise, in case of a circular orbit, x is undetermined.

1.10 Orbital Elements Defined for Any Orbit 91
When X and x are both determined, then the mean anomaly (M0) results from
M0 ¼ l0 X x
These orbital elements are due to Lagrange, who however used tan i instead of
tan(i/2), which choice leads to difficulties when i is near 90°. In 1774, Lagrange was
studying the secular effects produced by mutual planetary perturbations. The
advantage coming from using these elements resides in their application to orbits of
small eccentricity and/or small inclination, when the differential equations, which
describe the variations of the classical orbital elements due to the presence of
perturbing accelerations, become singular.
This type of singularity may also arise with the equinoctial elements, when the
orbit inclination angle (i) is equal to 180°. This singularity can be eliminated by
using either of two sets of equinoctial elements: the prograde (or direct) set and the
retrograde set. Thus, a broader definition of the equinoctial elements, which
incorporates the two sets, is given below [29]:
a q1 ¼ ½tanði=2ÞI sin X
p1 ¼ e sinðX þ IxÞ q2 ¼ ½tanði=2ÞI cos X
p2 ¼ e cosðX þ IxÞ l0 ¼ M0 þ X þ Ix
In these expressions, the quantity I, called the retrograde factor, is equal to +1 for
the direct set, or to −1 for the retrograde set of equinoctial elements. This means
that −x is used instead of x, and cot(i/2) instead of tan(i/2), when the retrograde set
is used [30]. As will be shown in Chap. 3, in the numerical integration of the
differential equations relating to the orbit of a satellite in the presence of pertur-
bations, it is necessary to convert, at each integration step, the equinoctial elements
into the components of the position (r) and velocity (r′) vectors with respect to an
inertial reference system, because the perturbing accelerations depend in general on
r and r′.
The modified equinoctial orbital elements, proposed by Walker et al. [31] are
another set of elements which may be used instead of those indicated above. The
modified equinoctial orbital elements are defined as follows
p ¼ að 1 e 2 Þ h ¼ tanði=2Þ cos X
f ¼ e cosðX þ xÞ k ¼ tanði=2Þ sin X
g ¼ e sinðX þ xÞ ‘0 ¼ X þ x þ /0
where p is the semi-latus rectum and ‘0 is the true longitude at epoch (see Sect. 1.9).
The choice of the semi-latus rectum instead of the major semi-axis makes it possible
to apply this set of elements to all orbits. The singularity arising in the h and
k elements when the orbit inclination angle (i) is equal to 180° can be removed
by using the retrograde set of modified equinoctial elements, as shown above.
Eagle [32] has shown how to compute the classical elements from the modified
equinoctial elements. Eagle’s formulae are given below
p 1
a¼ e ¼ f 2 þ g2 2
1 f2 g2
h 1 i h 1 i
i ¼ 2tan1 h2 þ k 2 2 ¼ tan1 2 h2 þ k2 2 ; 1 h2 k 2

g k
x ¼ tan1 tan1 ¼ tan1 ðgh fk; fh þ g k Þ
f h

g
X ¼ tan1 ðk; hÞ /0 ¼ ‘0 ðX þ XÞ ¼ ‘0 tan1
f
u0 ¼ x þ /0 ¼ tan1 ðh sin ‘0 k cos ‘0 ; h cos ‘0 þ k sin ‘0 Þ
where the expression tan−1(a, b) indicates a four quadrant inverse tangent calcu-
lation. Eagle has also shown how to compute the position vector (r) and the velocity
vector (r′) of the orbiting body in an Earth-centred inertial system X, Y, and Z by
using the modified equinoctial orbital elements. These formulae are given below
hr i hr
r¼ cos ‘ 0 þ a 2
cos ‘ 0 þ 2hk sin ‘ 0 uX þ ðsin ‘0 a2 sin ‘0
s2 s

2
2r
þ 2hk cos ‘0 ÞuY þ 2 ðh sin ‘0 k cos ‘0 Þ uZ
s
" 1 #
0 1 l 2
r ¼ 2 sin ‘0 þ a sin ‘0 2 h k cos ‘0 þ g 2 f h k þ a g uX
2 2
s p
" 1 #
1 l 2
þ 2 cos ‘0 þ a2 cos ‘0 þ 2 h k sin ‘0 f þ 2 g h k þ a2 f uY
s p
" 1 #
2 l 2
þ 2 ðh cos ‘0 þ k sin ‘0 þ f h þ g kÞ uZ
s p
where a2 = h2 − k2, s2 = 1 + h2 + k2, r = p/w, w = 1 + f cos ‘0 + g sin ‘0, and l is

the gravitational parameter of the Earth.
1.11 The Lagrangian Coefficients f, g, f′, and g′ in Closed

Form
As has been shown in Sect. 1.1, the motion of an orbiting body with respect to its
centre of attraction is confined to an invariable plane. This fact results from the
constancy of the angular momentum per unit mass (h). The plane of motion is
1.11 The Lagrangian Coefficients f, g, f′, and g′ in Closed Form 93
identified by the position (r) and velocity (r′) vectors of the orbiting body with
respect to a reference system having its origin in the centre of force. Unless a
trajectory is rectilinear, these two vectors are linearly independent of each other
and, as such, form a basis in the plane of motion.
In other words, given the position (r0) and velocity (r′0) vectors of the orbiting
body at an initial time t0, the position (r) and velocity (r′) vectors of the same body
at any subsequent time t can be expressed as a linear combination of r0 and r′0, as
follows
r ¼ f r0 þ g r00
r0 ¼ f 0 r0 þ g0 r00
where f, g, f′, and g′ are the coefficients of this linear combination. When these
coefficients are known as functions of time, then we can predict the future position
and velocity of the orbiting body.
Section 1.9 has shown how to express the Cartesian co-ordinates in the perifocal
reference system Oxyz, having its origin in the centre O of force and having also ux,
uy, and uz as its unit vectors along, respectively, x, y, and z, as follows

12
0 l

r ¼ ðsin /Þux þ ðe þ cos /Þuy
p
where r = a(1 − e2)/(1 + e cos /) and p = a(1 − e2).

Thus, at a given time t0, there results
r0 ¼ x0 ux þ y0 uy ¼ ðr0 cos /0 Þux þ ðr0 sin /0 Þuy

12
l

r00 ¼ x00 ux þ y00 uy ¼ ðsin /0 Þux þ ðe þ cos /0 Þuy
p
where r0 = a(1 − e2)/(1 + e cos /0) and p = a(1 − e2).

The expressions for f, g, f′, and g′ can be obtained, following Bate et al. [10] or
Campbell [33], by considering the vector products of r0 and r′0 by
r ¼ f r0 þ g r00
r0 ¼ f 0 r0 þ g0 r00
As to the Lagrangian coefficient f, the vector product r r′0 yields

r r00 ¼ f r0 þ g r00 r00 ¼ f ðr0 rÞ þ g r00 r00 ¼ f h ¼ f h uz
This is because the vector product of any vector by itself is the zero vector.
In addition, since

r r00 ¼ xy00 yx00 uz
then

f h uz ¼ xy00 yx00 uz
that is,
xy00 yx00
f ¼
h
Now, remembering that h = (lp)½ and the following expressions of Sect. 1.9
12
l
x ¼ r cos / x00 ¼ sin /0
p
12
l
y ¼ r sin/ y00 ¼ ðe þ cos /0 Þ
p
and inserting them into f = (x y′0 − y x′0)/h, there results after simplification
r
f ¼ ½e cos / þ cosð/ /0 Þ
p
By subtracting and adding unity in the expression between square brackets, there
results
r
f ¼ ½1 þ 1 þ e cos / þ cosð/ /0 Þ
p
Remembering that r = p/(1 + e cos /), there results 1 + e cos / = p/r, which in
turn, inserted into the preceding expression, yields
rh p i r r
f ¼ 1 þ þ cosð/ /0 Þ ¼ þ 1 þ cosð/ /0 Þ
p r p p
r r
¼ 1 ½1 cosð/ /0 Þ ¼ 1 ½1 cosðD/Þ
p p
where D/ = / − /0.
As to the Lagrangian coefficient g, the vector product r0 r yields

r0 r ¼ r0 f r0 þ g r00 ¼ f ðr0 r0 Þ þ g r0 r00 ¼ gh ¼ g h uz
In addition, since
r0 r ¼ ðx0 y y0 xÞuz
then
ghuz ¼ ðx0 y y0 xÞuz
that is,
x0 y y0 x
g¼
h
Now, remembering h = (lp)½ and the following expressions of Sect. 1.9
x0 ¼ r0 cos /0 x ¼ r cos /
y0 ¼ r0 sin /0 y ¼ r sin /
and inserting them into g = (x0y − y0x)/h, there results
1 rr0
g¼ 1 ½ðr0 cos /0 Þðr sin /Þ ðr0 sin /0 Þðr cos /Þ ¼ 1 sinð/ /0 Þ
ðl pÞ 2 ðl pÞ2
rr0
¼ 1 sinðD/Þ
ðl pÞ2
where D/ = / − /0.
As to the Lagrangian coefficient g′, the vector product r0 r′ yields

r0 f 0 r0 þ g0 r00 ¼ f 0 ðr0 r0 Þ þ g0 r0 r00 ¼ g0 h ¼ g0 h uz
In addition, since
r0 r0 ¼ ðx0 y0 y0 x0 Þuz
then
g0 h uz ¼ ðx0 y0 y0 x0 Þuz
that is,
x0 y0 y0 x0
g0 ¼
h
Now, remembering h = (lp)½ and the following expressions of Sect. 1.9

12
l
x0 ¼ r0 cos /0 x0 ¼ sin /
p
12
0 l
y0 ¼ r0 sin /0 y ¼ ðe þ cos /Þ
p
and inserting them into g′ = (x0y′ − y0x′)/h, there results after simplification
r0 r0
g0 ¼ ½e cos /0 þ cosð/ /0 Þ ¼ ½e cos /0 þ cosðD/Þ
p p
where D/ = / − /0. By subtracting and adding unity in the expression between

square brackets, there results
r0
g0 ¼ ½1 þ 1 þ e cos /0 þ cosðD/Þ
p
Remembering that r = p/(1 + e cos /), there results 1 + e cos /0 = p/r0, which
in turn, inserted into the preceding expression, yields

r0 p r0 r0 r0
g0 ¼ 1 þ þ cosðD/Þ ¼ þ 1 þ cosðD/Þ ¼ 1 ½1 cosðD/Þ
p r0 p p p
Finally, as to the Lagrangian coefficient f′, the vector product r′ r′0 yields

r0 r00 ¼ f 0 r0 þ g0 r00 r00 ¼ f 0 r0 r00 þ g0 r00 r00 ¼ f 0 h ¼ f 0 h uz
In addition, since

r0 r00 ¼ x0 y00 y0 x00 uz
then

f 0 h uz ¼ x0 y00 y0 x00 uz
that is,
x0 y00 y0 x00
f0 ¼
h
Before going any further in this way of expressing f′ as a function of the true
anomaly, we use, to this end, a relation which makes it easy to compute f′ when f, g,
and g′ are known. This relation is
r0 r00 ¼ r r0 ¼ h
The expression written above follows from the constancy of the angular
momentum per unit mass (h) in the two-body problem. Thus

r r0 ¼ f r0 þ g r00 f 0 r0 þ g0 r00 ¼ ff 0 ðr0 r0 Þ þ r g0 r0 r00

þ g f 0 r00 r0 þ g g0 r00 r00
Now, the first and the fourth of the four addends on the right-hand side of the
expression written above vanish, because the vector product of any vector by itself
is the zero vector. Thus, the expression reduces to

r r0 ¼ f g0 r0 r00 þ g f 0 r00 r0
h ¼ r0 r00 ¼ r r0 ¼ constant
and also remembering the identity a b = −b a, there results
h ¼ r00 r0
This expression, inserted into

r r0 ¼ f g0 r0 r00 þ gf 0 r00 r0
yields
h ¼ fg0 h gf 0 h
that is,
f g0 g f 0 ¼ 1
which makes it possible to express f′ as follows
fg0 1
f0 ¼
g
where the expressions for f, g, and g′ are those derived above.

Inserting these expressions into f′ = (f g′ − 1)/g yields after simplification
" 1 #
0 l 2
1 cosð D/ Þ rr0
f ¼ r0 r þ ½1 cosðD/Þ
p rr0 sinðD/Þ p
Now, since the expression
1 cos a a
¼ tan
sin a 2
is a trigonometric identity, then f′ is also expressible as follows

12
l D/ 1 cosðD/Þ 1 1
f0 ¼ tan
p 2 p r r0
In summary, the expressions of the Lagrangian coefficients f, g, f′, and g′ as

functions of the true anomaly / are as follows
r
f ¼ 1 ½1 cosðD/Þ
p
rr0
g¼ 1 sinðD/Þ
ð l pÞ 2
12
l D/ 1 cosðD/Þ 1 1
f0 ¼ tan
p 2 p r r0
r
g0 ¼ 1 ½1 cosðD/Þ
0
p
where D/ = / − /0 and D//2 = (/ − /0)/2. The Lagrangian coefficients can also

be expressed as a function of the eccentric anomaly Æ, as will be shown below.
We use again the perifocal reference system Fxyz, having its origin F in the
centre of force and such that the x, y-plane is the plane of motion of the orbiting
body.
By inspection of the preceding figure, the perifocal co-ordinates x and y of the

orbiting body are
x ¼ acos Æ ae ¼ að cos Æ eÞ
b b
y ¼ y ¼ ða sin Æ Þ ¼ b sin Æ
a a
Now, since for an ellipse there results

b2 ¼ a2 ðaeÞ2 ¼ a2 1 e2
then
x ¼ a ðcos Æ eÞ
1
y ¼ a 1 e2 2 sin Æ
Differentiating the preceding expressions with respect to time, there results
x0 ¼ a Æ 0 sin Æ
1
y0 ¼ a 1 e2 2 Æ 0 cos Æ
Now, since
1 l12
Æ0 ¼
r a
as will be shown below, then
a l12
1
0 0 ðlaÞ2
x ¼ a Æ sin Æ ¼ sin Æ ¼ sin Æ
r a r
að1 e2 Þ2 l12
1 1
0
1
2 2 0 ½lað1 e2 Þ2
y ¼ a 1 e Æ cos Æ ¼ cos Æ ¼ cos Æ
r a r
The equality Æ′ = (1/r)(l/a)½ follows from Kepler’s equation (see Sect. 1.3)
3 12
a
l
which, differentiated with respect to time, yields

3 12
a
1¼ Æ 0 ð1 e cos Æ Þ
l
Substituting
r ¼ að1 e cosÆ Þ
(which is the equation of an ellipse in polar co-ordinates r and Æ) into Kepler’s

equation differentiated with respect to time, there results
3 12
a r
1¼ Æ0
l a
which, solved for Æ′, yields
1 l12
Æ0 ¼
r a
The expressions derived above, that is,

1 1
0 ðlaÞ2 0 ½lað1 e2 Þ2

x ¼ sin Æ y ¼ cos Æ
r r
substituted into the equation
xy00 yx00
f ¼
h
also derived above, yield

1 1 1
aðcos Æ eÞ½lað1 e2 Þ2 cos Æ 0 þ að1 e2 Þ2 sin Æ ðlaÞ2 sin Æ 0
f ¼
hr0
Since
1
1
h ¼ ðl pÞ2 ¼ la 1 e2 2
then the expression written above becomes

1 1 1
aðcos Æ eÞ½lað1 e2 Þ2 cos Æ 0 þ að1 e2 Þ2 sin Æ ðlaÞ2 sin Æ 0
f ¼ 1
½lað1 e2 Þ2 r0
The preceding expression can be simplified as follows

a a
f ¼ ½ðcos Æ Æ Þcos Æ 0 þ sin Æ sin Æ 0 ¼ ðcos Æ cos Æ 0 e cos Æ 0
r0 r0
a a
þ sin Æ sin Æ 0 Þ ¼ ½cosðÆ Æ 0 Þ e cos Æ 0 ¼ ½cosðDÆ Þ e cos Æ 0
r0 r0
where DÆ = Æ − Æ0. Remembering again the equation of an elliptic orbit
r ¼ að1 e cosÆ Þ
and solving the equation for e cos Æ, we have

r
e cos Æ ¼ 1
a
The preceding expression, substituted into

a
f ¼ ½cosðDÆ Þ e cos Æ 0
r0
yields
ah r0 i a
f ¼ cosðDÆ Þ 1 ¼ 1 ½1 cosðDÆ Þ
r0 a r0
The preceding equation expresses the Lagrangian coefficient f as a function of

the eccentric anomaly in case of an elliptic trajectory. By operating likewise, it is
possible to derive the expressions of the other Lagrangian coefficients.
These expressions are given below [11]:
a
f ¼1 ½1 cosðDÆ Þ
r0
12
að r 0 v 0 Þ a
g¼ ½1 cosðDÆ Þ þ r0 sinðDÆ Þ
l l
1
0 ðlaÞ2
f ¼ sinðDÆ Þ
rr0
a
g0 ¼ 1 ½1 cosðDÆ Þ
r
1.12 The Lagrangian Coefficients f and g in Time Series
The Lagrangian coefficients f and g in closed form can be replaced by their

expressions in time series. The position vector r is expanded in a Taylor series
around epoch, as follows
ðt t0 Þ2 00 X1
ðt t0 Þn ðnÞ
r ¼ r0 þ ðt t0 Þr00 þ r0 þ . . . ¼ r0
2! n¼0
n!
where r(n)
0 is the n
th
derivative of the position vector r, with respect to time,
computed at epoch t = t0. Since r(n) lies in the rr′-plane, then r(n) is expressible as
follows
rðnÞ ¼ fn r þ gn r0
where f0, f1, …, fn and g0, g1, …, gn are the coefficients of the Taylor series for the
f and g expansions. The equation of motion
l
r00 þ r¼0
r3
is differentiated by using the Lagrange fundamental invariants e, k, and w, so called

because they are independent of the selected co-ordinate system and form a closed
set with respect to the operation of time differentiation. They are:
l
e¼
r3
r 0 r r0
k¼ ¼ 2
r r
v2 r 0 r 0
w¼ 2¼ 2
r r
As a result of this definition, we have
de lr 0 r0
e0 ¼ 3 4 ¼ 3e ¼ 3ek
dt r r
0 0 00
dk r r þ r r 2ðr r0 Þr 0 v2
k0 ¼ ¼ 2 e 2k2 ¼ w e 2k2
dt r2 r3 r
0 d w 1 dðr0 r0 Þ 2rr 0 v2 2r0 00 2v2 r 0
w ¼ 2 4 ¼ 2 r 3 ¼ 2ek 2w k
dt r dt r r r
¼ 2kðe þ wÞ
In order to calculate the first, second, third, … derivative of r, we differentiate
r0 ¼ v
r00 ¼ er
1.12 The Lagrangian Coefficients f and g in Time Series 103
using the results obtained above, that is,
e0 ¼ 3ek
k0 ¼ w e 2k2
w0 ¼ 2kðe þ wÞ
By so doing, we get
r0 ¼ v
r00 ¼ er
rIII ¼ 3ekr
ev
rIV ¼ 15ek2 þ 3ew 2e2 r þ 6ekv

rV ¼ 105ek3 45ekw þ 30e2 k r þ 45ek2 þ 9ew 8e2 v
..
.
Observing that the (n + 1)th derivative of the position vector r is
rðn þ 1Þ ¼ fn0 r þ fn v þ g0n v þ gn v0
and that
l
v0 ¼ r00 ¼ r ¼ er
r3
then
rðn þ 1Þ ¼ ðfn0 egn Þr þ ðfn þ g0n Þv
Hence, the following recurrence relations hold
fn þ 1 ¼ fn0 e gn
gn þ 1 ¼ fn þ g0n
If n = 0, then there results
rð0Þ ¼ f0 r þ g0 v
and therefore
f0 ¼ 1
g0 ¼ 0
f1 ¼ f00 eg0 ¼ 0
g1 ¼ f0 þ g00 ¼ 1
f2 ¼ f10 eg1 ¼ e
g2 ¼ f1 þ g01 ¼ 0
f3 ¼ f20 eg2 ¼ e0 ¼ 3ek
g3 ¼ f2 þ g02 ¼ e
f4 ¼ f30 eg3 ¼ 3e0 k þ 3ek0 þ e2 ¼ e 2e 15k2 þ 3w
g4 ¼ f3 þ g03 ¼ 3ek þ 3ek ¼ 6ek
f5 ¼ f40 eg4 ¼ e 105k3 45kw þ 36ek eð6ekÞ ¼ e 105k3 45kw þ 30ek

g5 ¼ f4 þ g04 ¼ e 2e 15k2 þ 3w þ ð6ekÞ0 ¼ e 45k2 þ 9w 8e
..
.
These results are the same as those found above by differentiating the two
expressions r′ = v and r′′ = −er. A computation procedure to determine the series
coefficients (f0, f1, …, fn; g0, g1, …, gn) for the f and g functions is indicated below.
s ¼ ðt t 0 Þ
! !
X
1 n
s ðnÞ
X
1 n
s X
1 n
s fn X
1 n
s gn
r¼ r0 ¼ ðfn r þ gn vÞ0 ¼ r0 þ v0
n¼0
n! n¼0
n! n¼0
n! n¼0
n!
X
1 n
s
f ð r 0 ; v0 ; t Þ ¼ ðfn Þ0
n¼0
n!
X
1 n
s
gð r 0 ; v 0 ; t Þ ¼ ð gn Þ 0
n¼0
n!
The f and g series are expressed as follows
1 1 1
f ¼ 1 e0 s2 þ e0 k0 s3 þ 2e0 15k20 þ 3w0 s4 þ
2 2 24
1 1
g ¼ s e 0 s 3 þ e 0 k0 s 4 þ
6 4
where e0, k0, and w0 are, respectively, e, k, and w at t = t0.

1.13 Canonical Units 105
1.13 Canonical Units
Canonical units are units of measurement related to a particular orbit around a given
celestial body. The latter may be, for example, the Sun or the Earth. The mass of the
chosen celestial body is taken as the unit of mass, and the mean distance between
this body and another body revolving about the former is taken as the unit of length.
When the Sun is chosen as the central body of reference, then the astronomical
system of units is used, which was adopted by International Astronomical Union
(IAU) in 1976. In the astronomical system of units, the distance unit is the astro-
nomical unit (AU), which is approximately the mean distance between the Earth
and the Sun. This unit of length is commonly used to express the distances of
bodies within the Solar System and is defined as the distance from the centre of the
Sun at which a particle of negligible mass, in an unperturbed circular orbit, would
have an orbital period of 2p/k or 365.2568983 days, where k is the Gaussian
gravitational constant, which is exactly equal to 0.01720209895.
One AU is exactly equal to 149597870700 m. In the same system of units, the
unit of mass is the mass of the Sun (1.9891 1030 kg), which is about 332950 the
mass of the Earth (5.9742 1024 kg) and 1048 times the mass of Jupiter
(1.8986 1027 kg). The astronomical unit of time is the day, defined as 86400 SI
seconds, where SI is the International System of Units.
By the way, the Gaussian gravitational constant k (so called after Carl Friedrich
Gauss, who proposed it in his work Theoria motus [34]) is equal to the square root
of the Newtonian gravitational constant G, when G is expressed in the astronomical
system of units. Let us remember Kepler’s third law
3 12
a
T ¼ 2p
GM
where T is the sidereal year (i.e., the time taken by the Earth to revolve about the
Sun once with respect to the fixed stars) expressed in days (T 365.256363 days),
a is the major semi-axis expressed in astronomical units (a = 1 AU), and M is the
mass of the Sun expressed in solar masses (MS = 1).
1
Now, Kepler’s third law, solved for G2 , yields
k G2 ¼ 0:01720209895 AU2 day1 M1

1 3
S
When the Earth is chosen as the central reference body, the canonical distance
unit (DU) is the distance between the centre of mass of the Earth and a fictitious
body revolving about the Earth in a circular orbit which barely touches the equator,
as shown in the following figure.
By so doing, the radius vector (r) and the velocity (v) of the fictitious orbiting
body are, respectively,
12
lE
r ¼ RE v¼
RE
where RE is the radius of the Earth at the equator, and lE = GME is the Earth
gravitational parameter. In this case, 1 DU is defined as follows:
1 DU ¼ RE ¼ 6378:137 km
where the value of RE given above is taken from the official Earth Gravitational
Model EGM2008 [35].
The canonical time unit (TU) is defined by the gravitational parameter of the
central reference body. Since the model cited above also gives the following value
lE = GME = 398600.4418 km3/s2 of the Earth gravitational parameter, then the
canonical time unit (TU) is defined as follows
3 12
RE
1 TU ¼ ¼ 806:811 s
lE
Consequently, the Earth gravitational parameter, expressed in canonical units, is
DU3
lE ¼ GME ¼ 1
TU2
The distances between the Earth and celestial bodies outside the Solar System
are so great as to require other units of length to be expressed in a practical way.
For example, the nearest star (except the Sun) to the Earth is the a Centauri
system, which is 271000 AU away. One of such units of length is the parsec
(parallax of one arcsecond), defined as the distance at which one astronomical unit
1.13 Canonical Units 107
subtends an angle of one arcsecond (1″), as shown in the following figure (not to
scale, due to the courtesy of NASA/GSFC [36]). One parsec is equivalent to about
30.857 1012 km or about 206264.8 AU.
By the way, parallax is, in general terms, an apparent change in the position of
an object resulting from a change in the observer position. In astronomy, parallax is
the angle subtended at a celestial body, especially a star, by the radius of the Earth
orbit (or, in other words, half the angle through which a star appears to move as the
Earth moves from one side of the Sun to the other). Consequently, the parallax of a
nearby star is determined by measuring the position of that star with respect to more
distant stars behind it, and then measuring again the position of the same star six
months later, when the Earth is on the opposite side of its orbit. When the star
considered is close enough to the Earth, a small but measurable parallax is
observed, because the position of the star has changed with respect to the more
distant stars placed in the background. The apparent change of position (2p) or
parallax angle (p) is very small (less than one second of arc) even for the nearest
star, as will be shown below.
Thus, with reference to the preceding figure, there results p tan(p) = 1/d,
where a is the radius of the Earth orbit measured in astronomical units, and d is the
distance to the star.
The a Centauri system has a parallax angle of 0.76 s of arc, corresponding to a

distance d = 1/p = 1/0.76 1.3 parsecs. The Sirius system has a parallax angle of
0.38 s of arc, corresponding to a distance of 1/0.38 2.6 parsecs.
Another unit of length frequently found in scientific articles is the light year,
defined as the distance travelled by light in one Julian year (i.e., in 365.25 days of
86400 SI s) in vacuo. Since light travels at a speed of exactly
2.99792458 1010 cm/s, a light year is about
2:99792458 1010 365:25 86400 9:46 1017 cm ¼ 9:46 1012 km
Thus, the a Centauri system is 4.3 light years, and the Sirius system is 8.6 light
years away. One parsec is equal to about 3.2616 light years.
1.14 The n-Body Problem
The n-body problem consists in determining the motion of an isolated set of

n bodies, having masses m1, m2, …, mn, attracting one another with Newtonian
gravitational forces, in the general case, when n is greater than two. Given the
positions and velocities of the n bodies at an initial time t0, it is required to
determine the positions and velocities of such bodies at any subsequent time
t. Following Battin [11], let
2 3
xi
6 7
xi 4 yi 5
zi
2 3
x0i
6 7
yi 4 y0i 5
z0i
be, respectively, the position vector and the velocity vector of the ith body, con-
sidered as a particle of mass mi, with respect to an inertial (i.e., non-accelerated and
non-rotating) reference system xyz. Let

1
rij ¼ xj xi xj xi 2
be the distance between two particles i and j, whose position vectors are xi and xj
and whose masses are mi and mj, attracting each other with a force of magnitude
mi mj
G
rij2
1.14 The n-Body Problem 109
where G = 6.67 10−11 Nm2/kg2 is the universal gravitation constant. The force
acting on mi due to mj is directed along the unit vector (xj − xi)/rij, whereas the
force acting on mj due to mi is oppositely directed. Consequently, the total force fi
acting on the particle mi, because of the presence of the other n − 1 particles, is
X
n
xj xi
0
fi ¼ G mi mj
j¼1
rij3
where the prime sign on the summation symbol does not denote a derivative with
respect to time, but only indicates the omission of the term for which i = j from the
sum. Newton’s second law of motion states that
d 2 xi dy
f i ¼ mi mi i
dt2 dt
because yi x′i. Hence, the following n differential equations in vector form
d2 xi Xn
xj xi
0
¼ G mi mj
dt 2
j¼1
rij3
with the appropriate initial conditions describe mathematically the motion of the
n particles. Since the motion of each particle mi is described by the three compo-
nents (xi, yi and zi) of its position vector xi and the three components (x′i, y′i and z′i)
of its velocity vector yi, then six times n integrals are needed for a system of
n particles. For example, in case of three bodies (n = 3), 6 3 = 18 integrals are
needed, because the solution is given by the three components of the position vector
xi and the three components of the velocity vector yi of the particle mi
xi ¼ xi ðt; c1 ; c2 ; . . .; c18 Þ x0i ¼ x0i ðt; c1 ; c2 ; . . .; c18 Þ

yi ¼ yi ðt; c1 ; c2 ; . . .; c18 Þ y0i ¼ y0i ðt; c1 ; c2 ; . . .; c18 Þ
zi ¼ zi ðt; c1 ; c2 ; . . .; c18 Þ z0i ¼ z0i ðt; c1 ; c2 ; . . .; c18 Þ
as functions of time and of eighteen integration constants, for each of the three
particles. Since the solution depends on eighteen constants, it implies the deter-
mination of eighteen integrals of motion, that is, eighteen equations involving only
algebraic functions or integrals in the co-ordinates and velocities of the three bodies
and satisfying the differential equations of motion. When no external forces act on
the particles, ten of these integrals are provided by the conservation of
• total linear momentum of the system;
• total angular momentum of the system; and
• total energy of the system.
As to linear momentum, the vector expression
X
n
xj xi
0
fi ¼ G mi mj
j¼1
rij3
shows that the total force f1 + f2 + + fn acting on the system of n particles has a
zero resultant. This means that the following vector
d2
ðm1 x1 þ m2 x2 þ þ mn xn Þ
dt2
has a zero resultant and hence no force acts on the centre of mass O of the system,
which point is identified by its position vector
m1 x1 þ m2 x2 þ þ mn xn
xO ¼
m1 þ m2 þ þ mn
In other words, the linear momentum of the system is conserved, and the
position vector of the centre of mass is expressible as follows
xO ¼ c 1 t þ c 2
where c1 and c2 are the vector constants of integration. Hence, the centre of mass of
the system is in a state of rest or of uniform motion in a straight line.
As to angular momentum, the vector expression
X
n
xj xi
0
fi ¼ G mi mj
j¼1
rij3
shows that not only the total force f1 + f2 + + fn, but also the sum of all vector
moments xi fi (for i = 1, 2, …, n) has a zero resultant.
This means that the following vector has a zero resultant, that is,

d d x1 d x2 d xn
m 1 x1 þ m2 x2 þ þ mn xn ¼0
dt dt dt dt
By integrating this expression and writing yi in place of dxi/dt (i = 1, 2, …, n),

there results
m 1 x1 y 1 þ m 2 x2 y 2 þ þ m n xn y n ¼ c 3
which indicates the constancy, in magnitude and direction, of the angular

momentum vector c3. Hence, the system of n particles has an invariable plane,
which is orthogonal to c3 and contains the centre of mass O of the particles.
The potential of gravitation Vi at the point Pi of co-ordinates xi, yi, and zi is a

scalar quantity defined as follows
Xn
0 mj
Vi ¼ G
j¼1
rij
Since the potential of gravitation Vi depends only on the distance between Pi and
the other particles, then Vi is also independent of the choice of the co-ordinate axes.
The importance of the potential of gravitation follows from the fact that, in a
gravitational force field, the force of attraction can be expressed as the gradient of Vi
on a particle of unit mass located at Pi whose co-ordinates are xi, yi, and zi. This
means that
@V i
f Ti ¼ mi
@xi
where ∂Vi/∂xi is defined as the following 1 3 row vector

@V i @V i @V i @V i

@xi @xi @yi @zi
The function U defined as follows
1X n
U¼ mi V i
2 i¼1
is called the force function and is equal to the total work done by the gravitational
forces in gathering together the n particles from infinite mutual distances to the
given configuration. The potential energy of the system of n particles is conse-
quently equal to −U. If the force function is used instead of the potential of
gravitation, the force vector fi, acting on the particle of mass mi, is expressible as
follows
@U
f Ti ¼
@xi
where U is independent of the co-ordinates of any particular point. The force

function U depends on the components x1, y1, z1, x2, y2, z2, …, xn, yn, zn of the
position vectors x1, x2, …, xn. Each of these components, in turn, depends on time,
so that the total derivative of the force function is
dU @U dx1 @U dy1 @U dzn

¼ þ þ þ
dt @x1 dt @y1 dt @zn dt
Since, as shown above, ∂U/∂xi is a row vector and dxi/dt is a column vector,
then
@U d xi @U dxi @U dyi @U dzi

¼ þ þ
@xi dt @xi dt @yi dt @zi dt
Hence, remembering the equation fTi = ∂U/∂xi, the total derivative of the force
function may be written as follows
dU d x1 d x2 d xn
¼ f T1 þ f T2 þ þ f Tn ¼ f T1 y1 þ f T2 y2 þ þ f Tn yn
dt dt dt dt
This equation, written in vector form, is
dU
¼ f 1 y1 þ f 2 y2 þ þ f n yn
dt
By applying Newton’s second law of motion fi = mi(dyi/dt) to the expression

written above, there results
dU dy dy dy dT
¼ m1 1 y1 þ m2 2 y2 þ þ mn n yn ¼
dt dt dt dt dt
where

1 0 2 0 2 0 2 1 0 2 0 2 0 2
T ¼ m 1 x1 þ y1 þ z 1 þ m2 x2 þ y2 þ z2
2 2
1 h 2 i
þ þ mn x0n þ y0n þ z0n
2 2
2
is the kinetic energy of the system of n particles. Hence, the total energy of the
system is
T U ¼c
where c is a constant. No further integrals than these are available in the general
case for the n-body problem. The ten constants of integration are the nine com-
ponents of the three vectors c1, c2, and c3 and the scalar constant c. However, if the
motion is considered intrinsically, not all of the missing constants of integration
(corresponding to the missing integrals) are necessary, since two of them are not
essential. This is because:
• there is not an explicit dependency on time in the differential equations of
motion (in particular, the gravitational forces do not depend explicitly on time),
and consequently, the origin of time can be chosen arbitrarily; and
• the two axes lying in the invariable plane orthogonal to the constant vector c3
can be chosen arbitrarily by assigning a value to a parameter (e.g., to the angle h
between the x-axis and a reference direction).
Though there are no analytical solutions in the general case of an n-body
problem when n is equal to or greater than three, analytical solutions are possible in
particular cases. As is well known, there are five equilibrium points L1, L2, L3, L4,
and L5, called Lagrangian points, in the vicinity of two celestial bodies having
masses m1 and m2. The following figure, due to the courtesy of NASA [37] shows
the five Lagrangian points.
Lagrange discovered these points in 1772, while he was studying the restricted
three-body problem, so called because the third mass (m3) is assumed to be neg-
ligible with respect to the other two (m1 and m2). As Lagrange wrote in his “Essai
sur le problème des trois corps”, he was searching such solutions of the restricted
three-body problem as would allow the mutual distances of the bodies, or at least
the proportions of their mutual distances, to remain constant. He found that these
conditions could only be satisfied if the three bodies were placed either on the same
straight line, or at the vertices of an equilateral triangle [38].
In the first case, the light body m3 can be in either of three (L1, L2, and L3) points
along the straight line joining the centres of mass of the two large bodies m1 and m2.
In the second case, the light body can be in either of two (L4 and L5) points placed
at the vertex of an equilateral triangle having the two large bodies at the other
vertices. Point L5 is the mirror reflection of L4 about the straight line joining the
centres of mass of the two large bodies. Consequently, in each system of two large
orbiting bodies (e.g., Earth-Moon or Sun-Jupiter) there are five Lagrangian
points, two (L4 and L5) of which are points of stable equilibrium provided that one
(e.g., m2) of the two large masses m1 and m2 is less than 0.0385 times the total mass
m1 + m2, meaning by this that a small third body m3 placed in either L4 or L5 would
be retained in its place in spite of slight perturbations due to outside gravitational
forces; whereas the other three (L1, L2, and L3) are points of unstable equilibrium.
In the Earth-Moon system, L1 and L2 are on either side of the Moon along the
line joining the Earth with the Moon, so that L1 is in front of the Moon, L2 is behind
the Moon and L3 is on the opposite side of the Earth with respect to the Moon.
L4 and L5 are along the orbit of the Moon, so that L4 is 60° ahead and L5 is 60°
behind the Moon.
In the Sun-Jupiter system, which satisfies the stability criterion
m2
\0:0385
m1 þ m2
the Trojan asteroids are known to be confined in the vicinity of the L4 and L5 points
of Jupiter’s orbit around the Sun. The results given above can be found by con-
sidering an inertial reference system xyz having its origin O in the centre of mass of
the two large bodies of mass m1 and m2.
Following Fitzpatrick [39], with reference to the following figure, let us assume
the xy-plane of this reference system as coinciding with the invariable plane con-
taining the orbits of the two large bodies around their centre of mass.
In addition, let us assume for simplicity the orbits of the two large bodies to be
circular, r1 and r2 being the distances of, respectively, m1 and m2 from their centre
of mass. Such an assumption is valid, because the eccentricities of the ellipses
described by the planets revolving around the Sun are generally small, as has been
shown in Sect. 1.1.
Let the two large masses m1 and m2 be aligned along the x-axis of the inertial
reference system xyz at the time t = 0, and let r = r1 + r2 be the constant distance
between them. The constant angular velocity x and the distance ratio r1/r2 are
h m þ m i12 r1 m 2
1 2
x¼ G ¼
r3 r2 m 1
where the first expression is third Kepler’s law and the second expression follows
from placing the origin O of xyz in the centre of mass.
If we choose arbitrarily the unit of length such that r = r1 + r2 = 1 and the unit
of mass such that G(m1 + m2) = 1, then there results
r1 m1 ¼ r2 m2
r1 m1 ¼ ð1 r1 Þm2
r1 ðm1 þ m2 Þ ¼ m2
m2 Gm2 Gm2
r1 ¼ ¼ ¼ ¼ l2
m1 þ m2 Gðm1 þ m2 Þ 1
r 2 ¼ 1 r 1 ¼ 1 l2
l2 ¼ 1 r2
Consequently, at any time t, the position vectors of the two large masses m1 and
m2 in the inertial reference system xyz with origin in O are, respectively,
2 3 2 3
l2 cos ðxtÞ l1 cos ðxtÞ
x1 4 l2 sin ðxtÞ 5 x2 4 l1 sin ðxtÞ 5
0 0
Now, let
2 3
x3
x3 4 y 3 5
z3
be the position vector of the small mass m3 at the same time t.

The three scalar equations of motion of m3 in the inertial reference system
defined above are
x3 x1 x3 x2
x003 ¼ l1 l2
q31 q32
y3 y1 y3 y2
y003 ¼ l1 l2
q31 q32
z z
z003 ¼ l1 3 l2 3
3 3
q1 q2
where
h i12
q1 ¼ ðx3 x1 Þ2 þ ðy3 y1 Þ2 þ z23
h i12
q2 ¼ ðx3 x2 Þ2 þ ðy3 y2 Þ2 þ z23
Let us consider now the motion of m3 with respect to a non-inertial reference

system rotating with the same angular velocity x about an axis orthogonal to the
plane of motion of the two large masses and passing through their centre of mass.
With reference to the following figure, let n, η and f be the Cartesian axes of this
non-inertial reference system having its origin in O, chosen so that the two large
masses m1 and m2 are always on the n-axis, the η-axis is orthogonal to the n-axis
and lies in the plane of motion of m1 and m2, and the f-axis has the same direction
as the z-axis of the inertial reference system.
As a result of this choice of axes, the position vectors n1 and n2 of, respectively,
m1 and m2 have the following constant components in the non-inertial reference
system nηf:
2 3 2 3
l2 l1
n1 4 0 5 n2 4 0 5
0 0
The position vector n3 of m3 has the following components, which vary with
time, in the same non-inertial reference system nηf:
2 3
n3
n3 4 g3 5
f3
In order to write the equations of motion of m3 in the non-inertial reference

system nηf, it is necessary to consider that m3 is subject not only to the Newtonian
gravitational acceleration written above for xyz, but also to two fictitious acceler-
ations. These fictitious accelerations are
(a) the Coriolis acceleration −2x n3
(b) the centrifugal acceleration −x (x n3)
where
2
3
0
x405
x
is the angular velocity vector. Consequently, the equation of motion of m3 in the

non-inertial reference system nηf is in vector form
d2 n3 n n n n
n003 ¼ l1 3 3 1 l2 3 3 2 2x n03 x ðx n3 Þ
dt2 q1 q2
where
h i12 h i12
q1 ¼ ðn3 n1 Þ2 þ g23 þ f23 ¼ ðn3 þ l2 Þ2 þ g23 þ f23
h i12 h i12
q2 ¼ ðn3 n2 Þ2 þ g23 þ f23 ¼ ðn3 l1 Þ2 þ g23 þ f23
The preceding differential equation in vector form, projected, respectively, onto

the n-axis, η-axis and f-axis, yields the following three scalar equations
n 3 þ l2 n l
n003 2xg0 3 ¼ l1 l2 3 3 1 þ x2 n3
q31 q2
0 g g
g003 þ 2xn 3 ¼ l1 3 l2 3 þ x2 g3
3 3
q1 q2
f f
f003 ¼ l1 33 l2 33
q1 q2
which can also be written as follows
@U
n003 2xg0 3 ¼
@n3
@U
g003 þ 2xn0 3 ¼
@g3
@U
f003 ¼
@f3
where
l1 l2 1 2
U¼ x n3 þ g23
q1 q2 2
is the sum of the gravitational and centrifugal potentials. Now, we ask ourselves
whether there may be any points (called Lagrangian points) in the force field
produced by the two large masses m1 and m2 such that a light mass m3 would
remain at rest if it were placed in one of these points. This means that m3 would
maintain the same distances with respect to m1 and m2, that is, that m3 would be in a
fixed position with respect to the rotating reference system nηf, and would rotate
with respect to the inertial reference system xyz with the same angular velocity x as
that of the large masses m1 and m2. According to the definition given above, the
Lagrangian points are such that m3, if placed in one of them, would have an
acceleration and a velocity equal to zero in the rotating reference system nηf
n003 ¼ 0
0
n3 ¼ 0
or, which is the same, the gradient of the force function U would be zero in the
rotating reference system
@U @U @U
¼ ¼ ¼0
@n3 @g3 @f3
Since
l1 l2 1 2
U¼ x n3 þ g23
q1 q2 2
then the partial derivative of U with respect to f3 is
@U @U @q1 @U @q2
¼ þ
@f3 @q1 @f3 @q2 @f3
remembering that
h i12 h i12
q1 ¼ ðn3 þ l2 Þ2 þ g23 þ f23 q2 ¼ ðn3 l1 Þ2 þ g23 þ f23
we have
@U l1 @U l2
¼ ¼
@q1 q21 @q2 q22

@q1 1 f @q2 1 f
¼ 2f3 ¼ 3 ¼ 2f3 ¼ 3
@f3 2q1 q1 @f3 2q2 q2
and therefore

@U l1 l2
¼ þ f3
@f3 q31 q32
Since l1/q31 + l2/q32 is positive definite, then ∂U/∂f3 can be equal to zero only
where f3 is also equal to zero. This means that the Lagrangian points must be in the
nη-plane (i.e., f3 = 0). On the other hand, if f3 = 0, then
h i12 h i12
q1 ¼ ðn3 þ l2 Þ2 þ g23 þ f23 ¼ ðn3 þ l2 Þ2 þ g23
h i12 h i12
q2 ¼ ðn3 l1 Þ2 þ g23 þ f23 ¼ ðn3 l1 Þ2 þ g23
l1 q21 þ l2 q22 ¼ l1 n23 2n3 l1 l2 þ l1 l22 þ l1 g23 þ l2 n23 þ 2n3 l1 l2 þ l2 l21 þ l2 g23

¼ n23 þ g23 ðl1 þ l2 Þ þ l1 l2 ðl1 þ l2 Þ ¼ n23 þ g23 þ l1 l2
This is because, by hypothesis, G(m1 + m2) = 1 and then l1 + l2 = 1. Thus, the

expression of U written above
l1 l2 1 2
U¼ x n3 þ g23
q1 q2 2
becomes

l1 l2 1 l l l 1 1 1 1
U¼ x 12 þ 22 1 ¼ l1 þ xq21 l2 þ xq22
q1 q2 2 q q2 l2 q1 2 q2 2
1
1 1 1 1 1 1
þ xl1 l2 ¼ l1 þ q2 l2 þ q2 þ l1 l2
2 q1 2 1 q2 2 2 2
This is because, by hypothesis, G(m1 + m2) = 1 and r = r1 + r2 = 1; conse-

quently, the angular velocity x = [G(m1 + m2)/r3]½ is also equal to unity. Thus, the
further two conditions which must be satisfied by the Lagrangian points, that is,
@U @U @q1 @U @q2 @U @U @q1 @U @q2

¼ þ ¼0 ¼ þ ¼0
@n3 @q1 @n3 @q2 @n3 @g3 @q1 @g3 @q2 @g3
become, respectively, after simplification

l1 1 q31 ðn3 þ l2 Þ l2 1 q32 ðn3 l1 Þ
þ ¼0
q31 q32

l1 1 q31 g3 l2 1 q32 g3
þ ¼0
q31 q32
The second of the two preceding equations for the two unknowns n3 and η3 is
evidently satisfied if η3 = 0, that is, by the Lagrangian points L1, L2, and L3 on the
n-axis. As shown in the preceding figure, L1 is between the positions of the two
large masses m1 and m2; L2 is on the right side of m2; and i3 is on the left side of m1.
In the case of L1, there results
n3 ¼ q1 l2 ¼ q1 ð1 l1 Þ ¼ l1 q2
q1 ¼ 1 q2
Substituting these expressions into

l1 1 q31 ðn3 þ l2 Þ l2 1 q32 ðn3 l1 Þ
þ ¼0
q31 q32
yields after simplification

1 l2 q32 13 q22 q2 þ 1
¼
3 l1 ð1 q2 Þ3 q22 þ q2 þ 1
which makes it possible to compute q2 for any given l2. For example, if l2 = 0.1,
then l1 = 1 − l2 = 1 − 0.1 = 0.9 and then q2 0.29.
Likewise, in the case of L2, there results
n3 ¼ q1 l2 ¼ q1 ð1 l1 Þ ¼ l1 þ q2
q1 ¼ 1 þ q2

l1 1 q31 ðn3 þ l2 Þ l2 1 q32 ðn3 l1 Þ
þ ¼0
q31 q32

1 l2 q32 13 q22 þ q2 þ 1
¼
3 l1 ð1 þ q2 Þ2 1 q32
which makes it possible to compute q2 for any given l2.

In the example given above (l2 = 0.1, l1 = 0.9), there results q2 0.36.
Likewise, in the case of L3, there results
n3 ¼ l2 q1 ¼ 1 þ l1 q1 ¼ l1 q2
q2 ¼ 1 þ q1

l1 1 q31 ðn3 þ l2 Þ l2 1 q32 ðn3 l1 Þ
þ ¼0
q31 q32

l2 1 q31 ð1 þ q1 Þ2
¼ 3 2
l1 q1 q1 þ 3q1 þ 3
In the example given above (l2 = 0.1, l1 = 0.9), there results q1 0.94.
As to the Lagrangian points L4 and L5, which do not lie on the n-axis, the two
conditions indicated above
@U @U @q1 @U @q2 @U @U @q1 @U @q2

¼ þ ¼0 ¼ þ ¼0
@n3 @q1 @n3 @q2 @n3 @g3 @q1 @g3 @q2 @g3
are evidently satisfied if
@U @U
¼ ¼0
@q1 @q2

1 1 1 1 2 1
U ¼ l1 þ q21 l2 þ q2 þ l1 l2
q1 2 q2 2 2
there results

@U 1 @U 1
¼ l1 2 þ q1 ¼ l2 2 þ q2
@q1 q1 @q2 q2
Thus, the condition
@U @U
¼ ¼0
@q1 @q2
is satisfied if
q1 ¼ q2 ¼ 1
Remembering that
h i12 h i12
q1 ¼ ðn3 þ l2 Þ2 þ g23 q2 ¼ ðn3 l1 Þ2 þ g23
and that
l1 ¼ 1 l2
the condition q1 = q2 = 1 yields
ðn3 þ l2 Þ2 þ g23 ¼ ðn3 1 þ l2 Þ2 þ g23 ¼ 1
The equation written above is satisfied if

1
1 32
n3 ¼ l2 g3 ¼
2 2
which are the Cartesian co-ordinates of the Lagrangian points L4 and L5 in the
rotating reference system nηf.
1.15 The Halo Orbits
By this term, coined by Farquhar [40] in 1968, we mean a family of symmetric,

periodic, three-dimensional trajectories surrounding one of the three collinear points
(L1, L2, and L3) of the five Lagrangian points in the restricted circular three-body
problem, which has been considered in Sect. 1.14. Strictly speaking, such trajec-
tories are not orbits, because they do not lay on a plane, nor do they enclose a
region of space occupied by an attracting body. They are placed, as mentioned
above, in the vicinity of one of the three collinear Lagrangian points, which are
points of unstable equilibrium. Therefore, station-keeping manoeuvres are required
in order for a spacecraft to move along one of such trajectories for a desired period
of time. As a possible application for these orbits, Farquhar has suggested to insert a
relay satellite into a halo orbit in the vicinity of the Lagrangian point L2 in the
Earth-Moon system, as shown in the following figure, which is due to the courtesy
of NASA [41].
1.15 The Halo Orbits 123
A single satellite, if forced to follow a halo orbit, would provide continuous

communications coverage between the Earth and the far side of the Moon. By so
doing, it would be possible to continuously control an unmanned roving vehicle on
the far side of the Moon [41]. The principal results found by Farquhar for this
application are briefly described below.
The nonlinear equations of motion for the circular restricted three-body problem
have been considered in Sect. 1.14. Let x, y, and z be the co-ordinates of the satellite
in the rotating reference system illustrated in the preceding figure. Farquhar [42] has
shown that these equations can be linearised as follows
y00 þ 2x0 þ ðBL 1Þy ¼ 0

x00 2y0 ð2BL þ 1Þx ¼ 0
z00 þ BL z ¼ 0
where BL = 3.19042. In these equations, the following quantities have been set
equal to unity: the unperturbed mean Earth-Moon distance (a = 384748.91 km),
the mean angular rate of the Moon around the Earth (n = 2.661699489 10−6
rad/s), and the sum (ME + MM) of the masses of the Earth (ME) and of the Moon
(MM), where ME/MM = 81.30. Farquhar [41] notes that the distance (d2) between
the Lagrangian point L2 and the Moon is not constant, but the ratio (cL) of this
distance to the instantaneous Earth-Moon distance (R) is constant, so that
d2
cL ¼ ¼ 0:1678331476
R
These equations show that the motion along the z-axis is independent of the
motion in the xy-plane, and is simply harmonic, with an angular frequency equal to
the square root of BL. They also show that the two differential equations governing
the motion in the xy-plane are coupled. The system of equations
x00 2y0 ð2BL þ 1Þx ¼ 0

y00 þ 2x0 þ ðBL 1Þy ¼ 0
can be solved by means of Laplace transforms [43]. By using this method, it is

possible to show that the motion of the satellite in the xy-plane has two components,
one of which is oscillatory, and the other is divergent. Because of the divergent
component, L2 is a point of unstable equilibrium, and therefore the satellite requires
station-keeping in order to remain near the L2 point.
These equations neglect the nonlinear terms, the eccentricity of the orbit of the
Moon around the Earth, and the perturbations of the gravitational field of the Sun.
Farquhar [41] has found that, under certain conditions, it is possible for the satellite
to follow a periodic halo orbit about L2, as shown in the following figure, which is
due to the courtesy of NASA [41].
To this end, the initial conditions must be chosen in such a way as to excite only
the oscillatory component of the motion in the xy-plane. However, since the
equilibrium along halo orbits is unstable, then station-keeping is necessary in order
for a satellite to remain in a given halo orbit. In particular, Farquhar [41] has found
that a whole family of halo orbits exists in the vicinity of L2, and that their
equations in the xy-plane are

xn ¼ Ax sin xxy t

yn ¼ Ay cos xxy t
where xxy = 1.86265 and Ax = 0.343336 Ay. The periods of these halo orbits are
about 14.67 days.
Again, the halo orbits, whose equations are given above, have been obtained
from the linearised differential equations of motion. Therefore, the halo orbits are a
1.15 The Halo Orbits 125
first-order approximation to the actual trajectories around the L2 point. The actual
trajectories can be obtained by finding particular solutions to the differential
equations of motion, which include nonlinear terms, the eccentricity of the orbit of
the Moon, and the perturbed gravitational field of the Sun.
Farquhar [41] has considered two cases for the actual trajectories:
(a) Lissajous paths, which are so called because the projection of each path onto
the yz-plane is a Lissajous curve. These trajectories are not periodic, because
they have a period of 14.65 days for the oscillatory motion in the orbital plane
of the Earth around the Moon, and a period of 15.23 days for the oscillatory
motion out of this plane.
(b) Halo paths, which are periodic trajectories, because the fundamental periods of
oscillation along, respectively, the y-axis and the z-axis are equal. This occurs
for every value of Ay greater than 32871 km.
An example of a Lissajous path is the trajectory of the WMAP (which stands for
Wilkinson Microwave Anisotropy Probe) around the Sun-Earth L2 point, which is
1.5 million km away from the Earth, as shown in the following figure, due to the
courtesy of NASA [44].
The lower portion of the preceding figure shows that Lissajous paths are inclined
with respect to the ecliptic plane, which is the plane containing the orbit of the
secondary body around the first, for the Sun-Earth system.
By the way, the WMAP spacecraft was launched on the 30th of June 2001. The
trajectory necessary to reach the observing station consisted of three lunar phasing
loops followed by an about 100-day cruise to the Sun-Earth L2 point. During this
cruise, a swing-by with the Moon was used to assist the spacecraft in reaching L2.
The spacecraft maintained a Lissajous path around L2, such that the WMAP-Earth
vector remained between 1° and 10° off the Sun-Earth vector to satisfy commu-
nications requirements while avoiding eclipses. Station-keeping manoeuvres were
required about four times per year to maintain the desired Lissajous path [44].
Finally, according to Kolemen et al. [45], there are three distinct families of
periodic orbits around the Lagrangian point L2: the horizontal Lyapunov orbits
(which are in the ecliptic plane), the horizontally-symmetric figure-eight-shaped
vertical Lyapunov orbits, and the halo orbits (which bifurcate from the family of the
horizontal Lyapunov orbits at higher values of energy). The orbits of the Lissajous
family cited above are not periodic.
References
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0-471-45703-5
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28(1A), 87–97 (2003). The cited article is also available at the web site http://www.kfupm.
edu.sa/publications/ajse/articles/281A_08P.pdf
20. R.J. Baun, P. Papadopoulos, Interplanetary mission design using multi-disciplinary design
optimization, article available at the web site http://ippw.jpl.nasa.gov/20070607_doc/4_
12BAUN.pdf
21. H.D. Curtis, Orbital Mechanics for Engineering Students (Butterworth-Heinemann, Oxford,
2005). ISBN 0-7506-6169-0
22. http://www.triumf.ca/project_ETA/gismo/gheisha_f77/rtmi.f
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24. http://celestrak.com/software/vallado-sw.asp
25. http://spaceflight.nasa.gov/realdata/elements/graphs.html
26. O. Montenbruck, E. Gill, Satellite Orbits (Springer, Berlin, 2005). ISBN 3-540-67280-X
27. “Eulerangles” by Lionel Brits—Hand drawn in Inkscape by me. Licensed under GFDL via
Wikimedia Commons—http://commons.wikimedia.org/wiki/File:Eulerangles.svg#mediaviewer/
File:Eulerangles.svg
28. J.L. Arsenault, K.C. Ford, P.E. Koskela, Orbit determination using analytic partial derivatives
of perturbed motion. AIAA J. 8(1), 4–12 (1970)
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(Mathematics Department, Naval Postgraduate School, Monterey, CA 93943). Web site
http://faculty.nps.edu/dad/orbital/th0.pdf
30. G.R. Hintz, Survey of orbit element sets. J. Guid. Control Dyn. 31(3), 785–790 (2008)
31. M.J.H. Walker, B. Ireland, J. Owens, A set of modified equinoctial orbit elements. Celest.
Mech. 36, 409–419 (1985)
32. C.D. Eagle, in Modified Equinoctial Orbital Elements, 2 June 2013, p. 9. Web site http://
www.cdeagle.com/pdf/mee.pdf
33. S. Campbell, in f and g Expressions (Coastal Bend College, Beeville Campus, Texas, United
States of America). article also available at the web site http://www.coastalbend.edu/Acdem/
math/campbel7/Stars/docs/f%20and%20g.pdf
34. C.F. Gauss, in Theoria Motus Corporum Coelestium in Sectionibus Conicis Solem
Ambientium (Perthes & Besser, Hamburg, 1809). Web site http://ia801409.us.archive.org/3/
items/bub_gb_ORUOAAAAQAAJ/bub_gb_ORUOAAAAQAAJ.pdf
35. NIMA Technical Report TR8350.2, Third edition, 3 Jan 2000. Web site http://earth-info.nga.
mil/GandG/publications/tr8350.2/wgs84fin.pdf
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gov/YBA/HTCas-size/parallax1-more.html
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38. J.L. Lagrange, in Essai sur le problème des trois corps, Oeuvres complètes, tome VI, pp. 229–331.
Web site http://sites.mathdoc.fr/cgi-bin/oeitem?id=OE_LAGRANGE__6_229_0
39. R. Fitzpatrick, in Lecture Notes on Analytical Classical Dynamics (University of Texas at
Austin). Web site http://farside.ph.utexas.edu/teaching/336k/lectures.pdf
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(Stanford University, Stanford, California, USA, NASA TR R-347). Web site https://ntrs.
nasa.gov/archive/nasa/casi.ntrs.nasa.gov/19710000821.pdf
41. R.W. Farquhar, in The Utilization of Halo Orbits in Advanced Lunar Operations, NASA TN
D-6365, July 1971, 101 p. Available at the web site https://ntrs.nasa.gov/archive/nasa/casi.
ntrs.nasa.gov/19710021579.pdf
42. R.W. Farquhar, in The Control and Use of Libration-Point Satellites, NASA TR R-346, Sept
1970, 125 p. Available at the web site https://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.gov/
19710000821.pdf
43. W.T. Thomson, in Laplace Transformation, 2nd edn. (Prentice-Hall, Englewood Cliffs, New
Jersey, USA, 1960)
44. NASA, WMAP observatory: trajectory and orbit, web site https://map.gsfc.nasa.gov/mission/
observatory_orbit.html
45. E. Kolemen, N.J. Kasdin, P. Gurfil, in Quasi-Periodic Orbits of the Restricted Three-Body
Problem Made Easy. New Trends in Astrodynamics and Applications III, American Institute
of Physics Conference Series, February 2007, vol. 886, Issue 1, pp. 68–77
Chapter 2
Orbit Determination from Observations
In general terms, the determination of an orbit is an iterative process meant to know

and predict the position and velocity (or the orbital elements) of a space object, with
respect to a primary celestial body, from observations of that object. For example,
the orbit determination for an artificial satellite revolving about the Earth is a series
of operations aimed at determining the motion (that is, the six orbital elements or
the six Cartesian components of the position and velocity vectors at some given
epoch) of the satellite with respect to a reference system having its origin in the
centre of mass of the Earth. Likewise, for a natural celestial body revolving about
the Sun, an astronomer computes the orbital elements of that body with respect to
the Sun on the basis of observations which have been performed at some place on
the surface of the Earth.
The methods used to this end may be classified into two broad categories. The
first category includes the classical (or deterministic) methods, which consider the
measurements as free from errors, and use therefore the minimum number of
measurements required to determine an orbit. The second category includes the
modern (or statistical) methods, which consider the measurements as affected by
errors, and use therefore more measurements than those which would be strictly
necessary to determine an orbit, with the view of reducing the influence of such
errors by means of a suitable mathematical treatment of the data gathered.
Let x = x(t) be the state vector, at a given time t, of an artificial satellite orbiting
around the Earth, that is, the vector whose six components are the three components
of the position vector, r, and the three components of the velocity vector, v, of the
satellite at time t with respect to either the true-of-date or the J2000.0
geocentric-equatorial reference system XYZ (defined in Sect. 1.9). Let x0 = x(t0) be
the known state vector of the same satellite at a given initial time t0 with respect to
the same reference system. When the forces (the central force and its perturbations)
acting upon the satellite are known exactly, then it is possible to determine its state

https://doi.org/10.1007/978-3-319-62220-0_2
130 2 Orbit Determination from Observations
vector x at any given time t by integrating the equations of motion. However, the
state vector x of the satellite at the initial time t0 is known approximately, not
exactly. In addition, the force model used to compute the forces acting upon the
satellite at a time t can only provide approximate values of the true forces, because
the physical constants (e.g. the gravitational parameter of the Earth) are known
approximately and also because the mathematical model used to compute the forces
can never be exact. Consequently, the orbit determination of a satellite requires that
the observations of the satellite (which are affected by random and systematic errors
due to the non-exactness of the force model used) should be mathematically treated
in such a way as to obtain the best estimate of the orbital elements representing the
satellite motion at any time. To this end, it was customary to operate in two stages.
The first stage, called preliminary orbit determination, is the approximate deter-
mination of an orbit by means of a minimum number of observations of the orbiting
object.
The second stage, called differential correction of orbits or orbit improvement,
comprises:
1. the collection of more observations than those strictly required, and
2. the fitting of the data gathered to an orbit by means of some mathematical
algorithm, based usually on the method of the least squares.
The results of the second stage are differential corrections to the preliminary
orbital elements, which have been determined in the first stage.
In the first stage, both the main attracting body and the orbiting body are con-
sidered as isolated particles subject only to their mutual gravitational forces.
Therefore, the motion is governed by the second-order differential equation shown
in Sect. 1.1, that is, by
l
r00 þ r¼0
r3
where r″ is the acceleration vector of the orbiting body, l is the gravitational

parameter of the attracting body (resulting from the product of the mass, M, of the
attracting body by the gravitational constant, G), and r is the position vector of the
orbiting body with respect to an inertial reference system having its origin in the
centre of mass of the attracting body, whose mass M is supposed to be much greater
than the mass m of the orbiting body. As also shown in Sect. 1.1, the constants of
integration of this differential equation are six, that is, the orbital elements of the
orbiting body.
In the second stage, the perturbations to the primary gravitational force are taken
into account, so that the differential equation of motion is expressed as follows
l
r00 þ r¼a
r3
where a is the vector sum of the perturbing accelerations acting on the orbiting
body. In case of an artificial satellite of the Earth, as will be shown at length in
Chap. 3, these accelerations are primarily due to:
• the non-spherical shape and the non-homogeneous mass density of the Earth;
• the gravitational attraction exerted on the satellite by other celestial bodies than
the Earth (in particular, by the Sun and the Moon);
• solid Earth and ocean tides;
• the aerodynamic drag exerted by the Earth atmosphere, particularly important in
case of artificial satellites orbiting at low altitudes;
• solar radiation pressure; and
• thrusters used for orbital manoeuvres.
By the way, solid Earth tides are similar to ocean tides, both of them being due
to the gravitational forces exerted upon the Earth by the Sun and the Moon. The
difference between the two types of tide resides in the much smaller tidal distortion
of the Earth (about 30 cm a day) in comparison with that of the ocean, because of
the resistance opposed by the rocks.
As a result of the perturbing accelerations, the orbital elements of an artificial
satellite revolving about the Earth are not constant, but vary slowly with time.
Consequently, the second-order differential equation r″ + (l/r3)r = a is
expressed as a system of m first-order differential equations
x0 ¼ f ðx; tÞ
with the initial conditions
xðt0 Þ ¼ x0
where x = x(t) is the augmented m-dimensional, time-dependent state vector, that is,
the state vector comprising not only the six position (x, y, z) and velocity (vx, vy, vz)
co-ordinates of the orbiting body but also the physical constants provided by the
force model used (that is, the physical quantities which are independent of time);
f(x, t) is a vector-valued function (the integrand function) which depends on the
state vector and on time; and x0 is the known state vector at some epoch t0.
On the other hand, the observations can be represented by the following system
of n nonlinear algebraic equations
zi ¼ gðxi ; ti Þ þ ei
( i = 1, 2, …, n), where zi are the actual observations made at epochs t1, t2, …, tn,
g(xi, ti) are the predicted values, and ei are random errors of measurement affecting
the observations. Such equations, solved for xi, yield the state vector
xðti Þ ¼ Hðx0 ; t0 ; ti Þ
(i = 1, 2, …, n) at times, respectively, t1, t2, …, tn. This way of proceeding,

comprising the two stages described above, is called batch estimation processing
and is based on the least-squares method.
Recently, as a result of the mathematical theory due to Kalman and Bucy [44,
45] and also of the advent of advanced computing means, the two stages are no
more separated from each other, and each epoch of observations is processed
individually by means of a sequential estimation algorithm, which is just the
Kalman filter. The preliminary orbit determination comprises several methods
meant to determine the orbital elements. These methods will be shown in the
following paragraphs. The determination of an orbit is part of a broader process
called integrated guidance, navigation, and control. The whole process of spacecraft
guidance, navigation, and control is shown in the following scheme, which is
redrawn from Thornton and Border [70].
Guidance is the actual steering of a spacecraft travelling through space. This

steering may come from sources placed either inside (a human crew or an on-board
computer) or outside (a ground station) the spacecraft. An example is provided by
the commands given to the rocket motor of a spacecraft to control the thrust.
Navigation is the measurement of the motion (position, orientation, and velocity) of
a spacecraft in space, obtained by means of observations performed by the crew, or
by automatic on-board sensors, or by tracking equipment located on the ground.
Control is the alignment and stabilisation of a spacecraft while the guidance and
navigation functions are being performed. These functions make it necessary not
only to determine the orbital elements (or the position and velocity vectors) of a
spacecraft, but also to adjust its path and attitude by means of control forces and
moments obtained by firing the main rocket motor or by routing commands to other
on-board devices (such as thrusters, reaction wheels, control moment gyroscopes,
or aerodynamic surfaces), which produce the desired forces on the spacecraft. The
functions of guidance and control are meant to compute and send a series of
commands to the propulsion system to alter the spacecraft velocity or attitude.
Following Thornton and Border [70], the orbit determination process requires an
a priori estimate (prediction) of the spacecraft trajectory, referred to as the nominal
orbit. In this estimate, expected values of the observable quantities are computed on
the basis of nominal values for the trajectory and models of such observable
quantities. The computed quantities differ from the correspondent observed quan-
tities coming from the tracking system. Such differences form the residuals (com-
puted minus observed quantities), which are due not only to random uncorrelated
measurement errors (such as thermal noise in the tracking receiver), but also to
errors in the trajectory and the observable models. The errors of the latter type
introduce distinctive signatures, that is, characteristic marks, in the residuals. Such
signatures make it possible to adjust the model parameters by means of a procedure,
called weighted linear least-squares estimation, which yields the set of parameter
values corresponding to the minimum weighted sum of the squares of the residuals.
This procedure, when the data are weighted by the inverse of their error
covariance, yields a minimum-variance estimator. Since this estimator provides a
linear solution to a nonlinear problem, then the estimation procedure comprises
more iterative steps, each of which uses the parameter estimation of the previous
step, until convergence is reached.
After the orbit determination process has been completed, the orbital elements
are compared with those required by the project. In case of discrepancies, trajectory
correction manoeuvres must be planned and executed. For example, when the
actual flight path differs from that which is required by the planned mission, it is
necessary to compute the magnitude and direction of the Dv vector required to
correct to the desired trajectory. In addition, suitable times must be computed, at
which the corrective manoeuvres will be executed. At the proper time, the space-
craft attitude will change the direction of the rocket motor axis to the desired one
and the thrusters will be fired for a determined interval of time. The magnitude of
Dv is usually small (metres or tens of metres per second), because of the limited
amount of propellant which can be carried on board.
Orbit trim manoeuvres are sometimes necessary in case of a satellite revolving
around the Earth, in order to prevent orbital decay due to air drag, and also to
perform orbit changes required by the mission objectives.
2.2 Topocentric Co-ordinate Systems
A radar station located on the surface of the Earth can measure the position and
velocity of an object, for example, of an artificial satellite, with respect to the place
where the station is located. However, the position vector r of interest to the
observer has its origin in another place, that is, in the centre of mass of the Earth. In
addition, the velocity vector v r′ of interest to the observer relates to the
geocentric-equatorial reference system XYZ, which moves with respect to the radar
station, because the Earth rotates about its axis. Thus, a series of measurements
relating to reference system located in the radar station must be converted so as to
provide the correspondent measurements with respect to the geocentric-equatorial
system. The measurements performed at a radar station are related to the
topocentric-horizon system UVW, which is defined below.
The topocentric-horizon reference system UVW, shown in the preceding figure,

has its origin in the point on the surface of the Earth where the radar station is
located, and its three Cartesian axes U, V, and W are directed so that its fundamental
plane UV is the horizon, with the U-axis pointing towards east, the V-axis pointing
towards north, and the W-axis perpendicular to the horizon and pointing upwards
(that is, towards the zenith). As to the directions forming the horizon, the agreement
of the authors is not unanimous. Some of them, for example Bate et al. [5] and
Vallado [77], choose the south and east directions. Other authors, for example
Montenbruck and Gill [53] and Curtis [20], choose the east and north directions.
We follow the latter convention.
2.2 Topocentric Co-ordinate Systems 135
Let the unit vectors of the three axes U, V, and W be, respectively, uU, uV, and
uW. Unlike the geocentric-equatorial reference system XYZ, which is fixed in space
with respect to the stars and does not rotate with the Earth, the topocentric-horizon
reference system UVW is evidently a non-inertial system, because it is fixed to the
Earth and does rotate with it.
Another reference system of the topocentric type is the topocentric-equatorial
co-ordinate system, shown on the left-hand side of the following figure.
The origin PH of this reference system is the point on the surface of the Earth
where the observer (or the radar station) is located. The axes Xt, Yt, and Zt of this
co-ordinate system are parallel to the axes, respectively, X, Y, and Z of the
geocentric-equatorial co-ordinate system having its origin O in the centre of the
Earth. The fundamental plane XtYt of the topocentric-equatorial system XtYtZt is
parallel to, but not coincident with, the fundamental plane XY of the
geocentric-equatorial system XYZ. Both X and Xt point towards the vernal equinox
♈. The position vector PHQ q, of magnitude q, goes from the point of obser-
vation PH to the space object Q. The position of Q is usually expressed by means of
the three co-ordinates q, d, and a, where the two angles d (declination) and a (right
ascension) are shown on the right-hand side of the preceding figure, which is an
enlarged view of the topocentric-equatorial co-ordinate system.
In the topocentric-equatorial system the North and South celestial poles are
determined by intersecting the rotation axis of the Earth with the celestial sphere,
that is, with the sky as seen from the Earth. The centre of the Earth is also the centre
of the celestial sphere, as shown in the following figure, which is due to the
courtesy of the University of Tennessee [76].
The celestial poles and equator are the geographic poles and equator projected
up onto the celestial sphere. The equivalent of the geographic latitude in the
topocentric-equatorial system is called declination and is measured in degrees North
(positive numbers) or South (negative numbers) of the celestial equator. The
equivalent of the geographic longitude in the topocentric-equatorial system is called
right ascension, and may also be measured in degrees, but for historical reasons it is
more common to measure it in time (hours, minutes, seconds), 1 hour of right
ascension being equivalent to 15 degrees of apparent sky rotation.
In addition to the celestial equator, another important plane intersecting the
celestial sphere is the ecliptic plane, which is the plane containing the orbit of the
Earth about the Sun, inclined of about 23°.4 with respect to the celestial equator, as
shown in the preceding figure. The inclination angle e 23°.4 is called the
obliquity of the ecliptic.
The position of a space object Q in the topocentric-equatorial system is
expressed by the following vector
q ¼ ðq cos d cos aÞuX þ ðq cos d sin aÞuY þ ðq sin dÞuZ
where uX, uY, and uZ are the unit vectors of, respectively, X, Y, and Z.
Let us express the position vector q as follows

q
q¼q ¼ quq
q
2.2 Topocentric Co-ordinate Systems 137
where q is the magnitude of q and uq = q/q is the unit vector having the direction
of q. By so doing, it is possible to write
uq ¼ ðcos d cos aÞuX þ ðcos d sin aÞuY þ ðsin dÞuZ
As shown above, the geocentric-equatorial system and the topocentric-equatorial

system have two different origins. Consequently, the direction cosines of the
position vectors r and q are not the same. Thus, the topocentric declination and
right ascension of a space object are not equal to the geocentric declination and
right ascension of the same object. However, when the magnitude of r is much
greater than the equatorial radius of the Earth, as is the case with stars and distant
planets, then the differences in topocentric and geocentric angles can be neglected.
2.3 Orbit Determination from a Single Radar Observation
With reference to the following figure, let us consider an artificial satellite inside the
field of view of a radar station located on the surface of the Earth.
The radar station measures the range, that is, the magnitude of the position
vector going from the radar station to the satellite and the direction of this position
vector with respect to a reference system having its origin in the place where the
radar station is located. Reference systems of this type are called topocentric,
because their origin is in the place of observation. One of these is the
topocentric-horizon reference system UVW, defined in Sect. 2.2 and also shown in
the preceding figure.
Let q be the position vector of an artificial satellite with respect to the
topocentric-horizon system indicated above. The direction of q is measured by two
angles, namely azimuth (A) and altitude (h), which result from the gimbal axes on
which the radar antenna is mounted.
The azimuth angle A is measured clockwise, for an observer in the direction of
the positive W-axis, from the north–south direction to the projection of q onto the
horizon (U, V) plane, so that 0 A 360°; the altitude angle h is measured
counterclockwise, for an observer in the direction of the positive U-axis, from the
projection of q onto the horizon plane to q itself, so that −90° h 90°.
If the station is equipped with a Doppler radar, that is, of a radar capable of
detecting a shift in frequency in the returning echo, then the rate of change of q can
also be measured. On the other hand, sensors on the gimbal axes can measure the
rate of change of the azimuth and altitude angles. Thus, the radar apparatus can
measure six quantities related to a satellite in its field of view, namely range (q),
azimuth angle (A), altitude angle (h), rate of change of range (q′), rate of change of
azimuth angle (A′), and rate of change of altitude angle (h′).
Let ‘U, ‘V, and ‘W be direction cosines of uq = q/q with respect to the
topocentric-horizon system UVW, which system has its origin in P and its axes
pointing to, respectively, east, north, and zenith. The unit vector uq = q/q is
expressible as follows
uq ¼ ‘U uU þ ‘V uV þ ‘W uW
where uU, uV, and uW are the unit vectors pointing to, respectively, east, north, and
zenith. By projecting uq onto, respectively, U, V, and W, there results
uq ¼ ðcos h sin AÞ uU þ ðcos h cos AÞ uV þ ðsin hÞ uW
Having obtained the components of the position vector q in the

topocentric-horizon reference system UVW, it remains to transform the same
components, for the purpose of computing the components of the position vector in
the geocentric-equatorial reference system XYZ. To this end, it is necessary to
determine the projections of the unit vectors uU, uV, and uW of the
topocentric-horizon system UVW onto the unit vectors uX, uY, and uZ of the
geocentric-equatorial system XYZ.
With reference to the following figure, let u be the geodetic latitude (defined
below) of the point P where the observer or the radar station is located, and let h be
the angular distance, measured in the equatorial plane, between the X-axis (pointing
towards the vernal equinox) and the local meridian passing through P.
2.3 Orbit Determination from a Single Radar Observation 139
In the preceding figure, the Earth is represented as an oblate ellipsoid, whose

equatorial bulge is exaggerated for the sake of clarity. The geodetic latitude u is the
angle between the equatorial plane XY and the normal in P to the surface of the
ellipsoid, which approximates the Earth. Since the Earth is not perfectly spherical,
then u does not coincide with the geocentric latitude, which is the angle between
the equatorial plane XY and the normal in P to the surface of the sphere.
Let uX* be the unit vector of OR, where OR lies in the plane of the meridian
passing through P and is perpendicular to the Z-axis. Let uW be the unit vector of
the zenith direction in P. By projecting uW onto OR and the Z-axis, there results
uW ¼ ðcos uÞ uX þ ðsin uÞuZ
where uZ is the unit vector of the Z-axis.

On the other hand, by projecting uX* onto X and Y, there results
uX ¼ ðcos hÞ uX þ ðsin hÞ uY
where uX and uY are the unit vectors of, respectively, X and Y.

It follows that
uW ¼ ½ðcos hÞuX þ ðsin hÞuY cos u þ ðsin uÞ uZ ¼ ðcos u cos hÞuX

þ ðcos u sin hÞuY þ ðsin uÞuZ
The unit vector uU, directed towards east, is expressed as a function of uX and
uY, by taking the vector product of uZ and uW, as follows
uZ uW
uU ¼
juZ uW j
where |uZ uW| is the magnitude of uZ uW. The following equality holds
2 3
uX uY uZ
a b ¼ 4 aX aY aZ 5
bX bY bZ
for any two vectors a = aX uX + aY uY + aZ uZ and b = bX uX + bY uY + bZ uZ.

2 3
uX uY uZ
6 7
uZ uW ¼ 4 0 0 1 5
cos u cos h cos u sin h sin u
¼ ð cos u sin hÞ uX þ ðcos u cos hÞ uY
h i12
juZ uW j ¼ ð cos u sin hÞ2 þ ðcos u cos hÞ2 ¼ cos u
uZ uW
uU ¼ ¼ ð sin hÞ uX þ ðcos hÞuY
juZ uW j
Finally, uV results from

2 3
uX uY uZ
6 7
uV ¼ uW uU ¼ 4 cos u cos h cos u sin h sin u 5
sin h cos h 0
¼ ð sin u cos hÞ uX ðsin u sin hÞ uY þ ðcos u cos2 h
þ cos u sin2 hÞ uZ ¼ ð sin u cos hÞ uX ðsin u sin hÞ uY þ ðcos uÞ uz
In summary, we have obtained the following results
uU ¼ ð sin hÞuX þ ðcos hÞuY

uV ¼ ð sin u cos hÞuX ðsin u sin hÞuY þ ðcos uÞuZ
uW ¼ ðcos u cos hÞuX þ ðcos u sin hÞuY þ ðsin uÞuZ
Let us consider the following scalar products
uU uX ¼ sin h
uU uY ¼ cos h
uU uZ ¼ 0
uV uX ¼ sin u cos h
uV uY ¼ sin u sin h
uV uZ ¼ cos u
uW uX ¼ cos u cos h
uW uY ¼ cos u sin h
uW uZ ¼ sin u
As shown in Sect. 1.9, in case of a transformation from perifocal co-ordinates

xyz to geocentric-equatorial co-ordinates XYZ, the rotation matrix R is defined as
follows
2 3
uX ux uX uy uX uz
R 4 uY ux uY uy uY uz 5
uZ ux uZ uy uZ uz
and this matrix is orthogonal, because it satisfies the condition RTR = I, where RT
is the transpose of R, and I is the 3 3 identity matrix.
Likewise, in the present case, the rotation matrix R, which transforms the
geocentric-equatorial co-ordinates XYZ into the topocentric-horizon co-ordinates
UVW, is defined as follows
2 3
uU uX uU uY uU uZ
R 4 uV uX uV uY uV uZ 5
uW uX uW uY uW uZ
Thus, by replacing the scalar products by their respective values which are given
above, there results
2 3 2 32 3
‘U sin h cos h 0 ‘X
4 ‘V 5 ¼ 4 sin u cos h sin u sin h cos u 54 ‘Y 5
‘W cos u cos h cos u sin h sin u ‘Z
where ‘U, ‘V, and ‘W are the direction cosines of the topocentric-horizon
co-ordinates, and ‘X, ‘Y, and ‘Z are the direction cosines of the
geocentric-equatorial co-ordinates.
The preceding expression defines the transformation from geocentric-equatorial
co-ordinates XYZ to topocentric-horizon co-ordinates UVW.
As is easy to verify, the 3 3 square matrix written above, that is,

2 3
sin h cos h 0
R 4 sin u cos h sin u sin h cos u 5
cos u cos h cos u sin h sin u
is an orthogonal matrix, because it satisfies the condition
RT R ¼ I
where RT is the transpose of R and I is the 3 3 identity matrix.

Therefore, in order to transform topocentric-horizon co-ordinates UVW into
geocentric-equatorial co-ordinates XYZ, we use RT, the transpose of R, that is,
2 3 2 32 3
‘X sin h sin u cos h cos u cos h ‘U
4 ‘Y 5 ¼ 4 cos h sin u sin h cos u sin h 54 ‘V 5
‘Z 0 cos u sin u ‘W
The matrices R and RT shown above can also be used for transformations from
topocentric-equatorial co-ordinates XtYtZt to topocentric-horizon co-ordinates UVW,
and from topocentric-horizon co-ordinates UVW to topocentric-equatorial
co-ordinates XtYtZt, because the unit vectors uX, uY, and uZ of the system XYZ
coincide with those of the system XtYtZt. The following relation holds between the
position vector (q) in the topocentric-horizon system UVW and the position vector
(r) in the geocentric-equatorial system XYZ:
r ¼ rE þ q
where rE OP is the vector going from the centre of mass, O, of the Earth to the
position, P, of the radar station on the surface of the Earth. This can be understood
by considering the following figure.
If the Earth were a perfect sphere having a radius equal to rE, then the local
vertical passing through the radar site P would join (extended downward) this site
with the centre O of the Earth. In this case, rE would be expressible as follows
rE ¼ rE u Z
Such is not the case in practice, because the shape of the Earth is not perfectly
spherical.
The following section of this paragraph takes account of the non-spherical shape
of the Earth. In the model described below, the Earth is approximated to an
ellipsoid, as shown in the following figure.
Consequently, latitude cannot be used any more as a spherical co-ordinate and

the radius of the Earth depends on latitude. Thus, a point on the surface of the Earth
will be represented by two rectangular co-ordinates (x and z) and the longitude (k),
because the latter has the same meaning in a spherical as in an oblate Earth. The
ellipsoid is an approximation to the true shape of the Earth.
The cross section of the Earth ellipsoid along a meridian is an ellipse, whose
major semi-axis (aE) is equal to the equatorial radius of the Earth and whose minor
semi-axis (bE) is equal to the polar radius of the Earth. Cross sections of the
ellipsoid parallel to the equator are circles. The shaded area represents the bulge,
which is maximum (21.4 km) at the equator and is zero at the poles.
The values considered here for the major semi-axis aE and the flattening
(or oblateness) f = (aE − bE)/aE of the Earth ellipsoid are those given in Ref. [58],
that is,
1
aE ¼ 6378:137 km f ¼ 0:0033528
298:257223563
Hence, the eccentricity (eE) and the minor semi-axis (bE) of the ellipsoid result,
respectively, from eE = (2f − f2)½ = 0.08181919 and bE = aE(1 − f) = 6356.752 km.
The reference ellipsoid defined above is an approximation to the geoid, that is, to the
equipotential surface of the Earth gravity field which best fits, in the least-squares
sense, the global mean sea level. In describing the non-spherical shape of the Earth,
some authors, Herrick [34, 35] for one, use f instead of eE.
The angle u* between the equatorial plane and the radius OP (between the
centre O of the Earth and the point P of the Earth surface where the observer is
located) is called geocentric latitude. The angle u between the equatorial plane and
the normal in P to the surface of the ellipsoid is called geodetic latitude and is used
in the maps and charts of the Earth. The normal in P to the surface of the ellipsoid is
the direction which a plumb bob would indicate, were it not for local anomalies in
the gravitational field of the Earth. The angle ua between the equatorial plane and
the actual plumb bob vertical uncorrected for these anomalies is called the astro-
nomical latitude. In practice, since the difference between the reference ellipsoid
and the mean sea level is small, then the difference between geodetic latitude u and
the astronomical latitude ua is negligible.
We must now calculate the station co-ordinates of a point on the surface of the
reference ellipsoid, when we are given the geodetic latitude, the longitude, and the
height of that point above the mean sea level, which is assumed to be the height
above the reference ellipsoid. To this end, with reference to the following figure,
we first determine the co-ordinates x and z of a point P on the ellipse, assuming that
we know the geodetic latitude u of P; then these co-ordinates will be adjusted for
another point of which we know the height above the surface of the ellipsoid in the
direction of the normal in P.
The angle b, called the reduced altitude, is introduced to express the x and z
co-ordinates as functions of the equatorial radius of the Earth and of b itself.
By considering the elliptic cross section of the Earth (having aE as its major
semi-axis along the equator and bE as its minor semi-axis along the polar axis of the
Earth) and the corresponding auxiliary circumference (having aE as its radius), the
Cartesian co-ordinates of P are expressible as follows

bE
x ¼ aE cos b z¼ aE sin b
aE
In Sect. 1.3, we have shown that, for an ellipse, there results

1
bE ¼ aE 1 e2E 2
where eE is the eccentricity of the ellipse. Thus, there also results

1
x ¼ aE cos b z ¼ aE 1 e2E 2 sin b
Now, cos b and sin b must be expressed as functions of the geodetic latitude u
and of the constant quantities aE and bE. To this end, since the slope of the tangent
to the ellipse is dz/dx, and the slope of the normal is −dx/dz = tan u, then the
differentials of the expressions written above
1
x ¼ aE cos b z ¼ aE 1 e2E 2 sin b
are
1
dx ¼ aE sin b db dz ¼ aE 1 e2E 2 cos b db
and consequently
dx tan b
tan u ¼ ¼
dz ð1 e2 Þ12
E
and also
1 1 sin u
tan b ¼ 1 e2E 2 tan u ¼ 1 e2E 2
cos u
Now, the last expression can be written in the following form
A
tan b ¼
B
where A = (1 − e2E)½ sin u and B = cos u. In addition, since

tan a 1
sin a ¼
1 cos a ¼
1
ð1 þ tan2 aÞ 2
ð1 þ tan2 aÞ2
are trigonometric identities, then there results after simplification

1
1 e2E 2 sin u cos u
sin b ¼
1 cos b ¼
1
1 e2E sin2 u 2 1 e2E sin2 u 2
The preceding expressions make it possible to write the x and z co-ordinates of a

point P on the surface of the ellipsoid as follows
aE cos u
x ¼ aE cos b ¼ 1
1 e2E sin2 u 2

1 aE 1 e2E sin u
z ¼ aE 1 eE sin b ¼
2 2
1
1 e2E sin2 u 2
Let us consider now another point PH, placed at a height H above the ellipsoid,
that is, at a height H above the mean sea level, along the normal in P to the
ellipsoid. By inspection of the preceding figure, the x and z components of the
height H result
Dx ¼ H cos u
Dz ¼ H sin u
The quantities Dx and Dy are the increments which must be added to the
co-ordinates of P derived above, to obtain the co-ordinates of PH as functions of
geodetic latitude (u), altitude (H) above the mean sea level, and equatorial radius
(aE) and eccentricity (eE) of the ellipsoid representing the Earth:
2 3 2 3
aE aE 1 e2E
xH ¼ 4 5
1 þ H cos u zH ¼ 4 5
1 þ H sin u
1 e2E 2
sin u 2 1 e2E sin2 u 2
When the flattening (f) of the Earth ellipsoid is used instead of eE, that is, when
e2E is replaced by 2f − f2, then the expressions derived above become
8 9
< aE =
xH ¼ þ H cos u
:1 ð2f f 2 Þ sin2 u 12 ;
8 9
< aE ð 1 f Þ 2 =
zH ¼ þ sin u
: 1 ð2f f 2 Þsin2 u 12
H
;
The third co-ordinate of PH is the east longitude (k) of PH. Thus, if we know the
Greenwich sidereal time (hG), it can be used with the east longitude of PH to
compute the local sidereal time (h), which is the angle between the X-axis (pointing
towards the vernal equinox) of the geocentric-equatorial system XYZ and the local
meridian. As will be shown below, the xH and zH co-ordinates plus the angle h, in
turn, are used to locate the observer site in the geocentric-equatorial system.
In the model representing the Earth as a perfect sphere of radius rE, rE was the
vector from the centre of the Earth to the position of the radar station on the surface
of the Earth, with respect to the geocentric-equatorial system XYZ.
Now, likewise, in the model representing the Earth as an oblate ellipsoid, let rE
be the position vector from the centre, O, of the Earth to the point PH, with respect
to the same system XYZ. The components of rE are then
rE ¼ ðxH cos hÞuX þ ðxH sin hÞuY þ zH uZ
where uX, uY, and uZ are the unit vectors along, respectively, X, Y, and Z.
Now, the vectors which have been expressed in the topocentric-horizon (or
UVW) co-ordinates must be converted into the geocentric-equatorial (or XYZ)
co-ordinates. The angle hG from the unit vector uX (pointing towards the vernal
equinox) and the Greenwich meridian is the Greenwich sidereal time. This angle,
added to the geographic longitude k of the radar site measured eastward from
Greenwich, yields the local sidereal time h, as follows
h ¼ hG þ k
The angles u and h determine the relation between the topocentric-horizon system
and the geocentric-equatorial system. To this end, it is necessary to compute hG at
any given time t. When hG is known, h results from the expression (h = hG + k)
given above. Let hG0 be the Greenwich sidereal time at some particular time t0 (e.g.
t0 = 0 h Universal Time on the 1st of January of the year of interest). When hG0 is
known, then the local sidereal time h at any given time t can be determined from
h ¼ hG 0 þ xE ðt t0 Þ þ k
where xE = 7.292 10−5 rad/s is the angular velocity of the Earth about its axis.
The following paragraph will show how to compute hG0, that is, hG at 0 h Universal
Time for every day of the year of interest. The same paragraph will also provide
further information on the concept of sidereal time introduced here.
It is to be noted that the vector rE OPH shown above is the position vector
from the centre, O, of the Earth to the point PH (where the observer is located), with
respect to the geocentric-equatorial system XYZ. The knowledge of rE and q, where
q PHQ is the position vector of a celestial body Q in the topocentric-horizon
system UVW, suffices for orbits of bodies revolving about the Earth, because
r = rE + q, where for such orbits r is the position vector from the centre of mass, O,
of the Earth to the body Q revolving in a geocentric orbit.
Following Boulet [12], we consider now bodies revolving about the Sun. With
reference to the following figure, let r be the position vector of a body B revolving
about the Sun, with respect to a celestial reference system XYZ whose origin is the
centre of mass of the Sun, and whose plane XY is the Earth equator.
Let E and S be the centres of mass of, respectively, the Earth and the Sun. Let q
be the topocentric position vector of the body B. As shown below, the knowledge
of rE EPH and s ES is necessary to compute PHS = −EPH + s.
The vector EPH defines the position of the observer with respect to the centre of
the Earth at time t. The local sidereal time h, the latitude u, and the east longitude k
are used to compute the rectangular components of EPH.
The most common unit of measurement for distances of celestial bodies within
the solar system is the astronomical unit (AU), which has been defined in Sect. 1.13
.
Let aE = 4.263523 10−5 AU (from Ref. [12]) be the radius of the Earth
expressed in astronomical units (AU). The rectangular components of rE EPH at
time t are
rEX ¼ aE cos u cos h

rEY ¼ aE cos u sin h
rEZ ¼ aE sin u
As to the components XS, YS, and ZS of the solar vector s ES, the geocentric
rectangular equatorial co-ordinates of the Sun, with respect to the mean equator and
equinox of J2000.0 (the 1st of January 2000 at noon), are published each year, for
every day of the year, in “The Astronomical Almanac”. The tabular values of the
co-ordinates of the Sun can be interpolated so as to obtain the values relating to the
time of interest. Alternatively, a simple algorithm can be used for computing the
angular co-ordinates of the Sun to an accuracy of about 1 arc-minute within two
centuries of 2000, which is given by the U.S. Naval Observatory [72]. This and
other similar algorithms require the knowledge of concepts on the measurement of
time in astronomy which have been not yet introduced to the reader and will be
shown at the end of Sect. 2.4.
When XS, YS, and ZS are known, the components of PHS result from
ðPH SÞX ¼ ðaE cos u cos hÞ þ XS

ðPH SÞY ¼ ðaE cos u sin hÞ þ YS
ðPH SÞZ ¼ ðaE sin uÞ þ ZS
and therefore r is given, in the heliocentric system XYZ, by
r ¼ PH S þ q
2.4 The Measurement of Time in Astronomy
As has been shown in the preceding paragraph, the determination of the orbit of
either an artificial satellite or a natural celestial body by means of observations
makes it necessary to record not only the observations themselves but also the times
at which they have been performed.
The time commonly used in every-day life is the time indicated by ordinary
clocks, that is, the mean solar time. The measurement of solar time is based on the
apparent motion of the Sun through the sky, due to the rotation of the Earth about
its axis, as seen by an observer placed in a fixed point of the surface of the Earth.
Thus, a solar day is the interval of time taken by the Sun, in its apparent motion
around the Earth, to travel an arc of 360° along the sky. In other words, a solar day
is the time elapsed between two consecutive passages of the Sun across the same
meridian of the place of observation. By meridian of the place of observation we
mean the great circle passing through the two celestial poles and the zenith of the
site. A solar day is measured by the astronomers from noon to noon, and is divided
into 24 h, or 24 60 = 1440 min, or 24 60 60 = 86400 s.
However, the rotation of the Earth about its axis does not comprise exactly 360°
in a solar day, because the Earth also revolves about the Sun. In order for the Earth
to do one revolution about the Sun, that is, in order for the Earth to travel an arc of
360° about the Sun, 365.25 solar days are necessary.
Thus, 360° travelled in 365.25 solar days are about 0°.985626 per solar day.
This means that, in the course of a solar day, the direction of the Sun seen from a
point of the Earth changes by about 1°. This also means that the Earth must travel
an arc of about 361° in order for the Sun to travel (in its apparent motion with
respect to the Earth) an arc of 360° along the sky.
The astronomers are interested to determine the rotation time of the Earth about
its axis not only with respect to the Sun, but also with respect to the distant stars (the
so-called fixed stars). In other words, they want to determine how long it takes the
Earth to rotate 360° around its axis with respect to the fixed stars. This period of
rotation is called sidereal day and can be determined by observing the starry sky by
2.4 The Measurement of Time in Astronomy 151
night. The difference between solar day and sidereal day is illustrated in the pre-
ceding figure (due to the courtesy of Wikimedia, Ref. [1]), where the amplitude of
the angle a = 0°.985626 has been exaggerated for the sake of clarity.
When the Earth is in position A along its orbit, the Sun and a given star of
reference are, both of them, overhead. In position B, the Earth has performed a
complete rotation about its axis, so that the star of reference (but not the Sun) is
overhead again. The time interval taken by the Earth to move from A to B equals
one sidereal day. Shortly after time tB, that is, at time tC, the Sun is overhead again.
The time interval taken by the Earth to move from A to C equals one solar day.
Thus, two consecutive passages of a fixed star (chosen as the star of reference)
across the same meridian of the Earth measure one sidereal day, which is also
divided into sidereal hours, minutes, and seconds.
Since two consecutive passages of the star of reference through the same
meridian take about 23 h, 56 min and 4 s of ordinary solar time to occur, then a
sidereal day is on the average 3 min and 56 s shorter than a solar day. In other
terms, a solar day is longer than a sidereal day by a factor which is about
360 þ 0:985626
¼ 1:00273785
360
The United States Naval Observatory [75] indicates a value of 1.00273790935

for this factor. This is the value which will be used below.
In general terms, sidereal time is defined by some authors (see, e.g., Refs. [34,
53]) as the hour angle of the vernal equinox, that is, the hour angle of the ascending
node of the ecliptic on the celestial equator. This is because the daily motion of the
vernal equinox measures the rotation of the Earth with respect to the fixed stars, not
to the Sun. By the way, the hour angle of any given point is the angle between the
half plane determined by the rotation axis of the Earth and the zenith (half of the
meridian plane) and the half plane determined by the rotation axis of the Earth and
the given point. This angle is taken with the minus sign if the given point is
eastward of the meridian plane, and with the plus sign if the given point is westward
of the meridian plane. The hour angle is usually expressed in time units (hours,
minutes, and seconds), where 24 h correspond to 360°. In particular, a sidereal day
is the interval of time taken by the hour angle of the vernal equinox to increase by
360°.
There are no two solar days having the same duration. This is because the axis of
rotation of the Earth is not perpendicular to the ecliptic (that is, to the plane
containing the orbit of the Earth around the Sun), and also because the orbit of the
Earth is slightly elliptical. Since the areal velocity of the Earth is constant, the Earth
moves faster along its orbit at perihelion (early in January) than it does at aphelion
(early in July). Thus, a mean solar day is defined by making reference to the Earth
revolving about the Sun in a circular orbit placed in the same plane as that of the
ecliptic and having the same period as that of the real elliptical orbit [5].
The difference between the true and the mean solar time is called the equation of
time. The two causes (that is, the eccentricity of the Earth orbit and the obliquity of
the ecliptic) mentioned above produce effects which overlap with different periods
of time, because the eccentricity has a period of one year, whereas the obliquity of
the ecliptic has a period of half a year. Consequently, the equation of time has two
minima and two maxima per year, as has been shown by Husfeld and Kronberg
[38].
As we all know, the Earth comprises 24 time zones, equally spaced in longitude
of about 15°, so that the mean solar time of each zone differs by ± 1 h from the
mean solar times of the two contiguous time zones. The world time zones are
shown in the following figure, due to the courtesy of CIA [16].
Of all these mean solar times, that which relates to the Greenwich meridian is
called Greenwich Mean Time (GMT) or Universal Time (UT1) or Zulu time (Z).
There are several versions of Universal Time, the principal of which are the
following:
• UT0 is the mean solar time of the Greenwich meridian obtained from direct
astronomical observations.
• UT1 is UT0 corrected for the effects of small movements of the Earth relative to
the axis of rotation (polar variation).
• UT2 is UT1 corrected for the effects of a small seasonal fluctuation in the rate of
rotation of the Earth. UT2 is mainly of historic interest and rarely used.
Universal Time is based on the imaginary mean Sun, which takes into account
the effects on the solar day of the weakly elliptic orbit of the Earth about the Sun.
As shown above, UT1 is a measure of the rotation angle of the Earth around its axis
as resulting from astronomical observations, account being taken of slight move-
ments of the Earth poles of rotation. UT1 predicts the solar position with sufficient
accuracy for astronomical purposes, but the duration of a second derived from UT1
varies noticeably because of variations of the Earth rotation (the velocity of the
Earth rotation is variable).
As has been shown by various authors (see, e.g., McCarthy [52]), these varia-
tions may be classified into three types: secular, irregular, and periodic.
The secular variation of the Earth rotational velocity depends on the apparently
linear increase in the length of the day due mainly to tidal friction. The Moon and
(to a lesser extent) the Sun raise tides in the oceans. Friction carries the maximum
tide ahead of the line joining the centres of the Earth and Moon. The resulting
couple decreases the velocity of rotation of the Earth and increases the orbital
momentum of the Moon. In other words, the Earth loses energy and slows down,
whereas the Moon gains this energy and increases its orbital period and distance
from the Earth. According to McCarthy, the decrease of the Earth rotational
velocity results in an increase of the day duration by about 0.0005 to 0.0035 s per
century. According to Espenak and Meeus [23], the secular acceleration of the
Moon implies an increase in the day length of about 0.0023 s per century. The
irregular variation of the Earth rotational velocity appears to be due to random
accelerations, but may be correlated with physical processes occurring on or within
the Earth. The resulting variation in the day length is evaluated by McCarthy to
0.001 s over the past 200 years.
Finally, the periodic variation of the Earth rotational velocity is associated with
periodically repeatable physical processes affecting the Earth. According to
McCarthy, tides raised in the solid Earth by the Moon and the Sun produce changes
in the length of the day with amplitudes of the order of 0.00005 s and with periods
of 18.6 years, 1 year, ½ year, 27.55 days, 13.66 days, and others. Due to the
reasons indicated above, the rotational velocity of the Earth varies in an unpre-
dictable manner.
In the most common civil usage, Universal Time is related to a time called
Co-ordinated Universal Time (UTC), which provides the basis for Civil Time.
UTC is kept by several time laboratories around the world and is measured by
high-precision atomic clocks. The length of a UTC second is defined in terms of an
atomic transition of Caesium under specific conditions and does not depend on the
observation of astronomical phenomena. The international standard UTC is pro-
vided by the International Bureau of Weights and Measures on the basis of the data
coming from the timing laboratories and is accurate to about 1 10−9 s (that is,
1 ns) per day. UTC is made public by various radio stations, which also provide the
difference between UTC and UT1, so that the latter can be computed as a function
of the former. As mentioned above, UTC is the basis to compute Civil Time in the
time zones of the Earth. By international agreements, UTC (which is measured by
atomic clocks) cannot differ from UT1 (which is measured by the rotation of the
Earth) by more than 0.9 s. When this limit is approached, a one-second change
(called leap second) is introduced into UTC. The UTC can be easily computed
starting from the local Civil Time (CT), which depends on where the observer is
placed, as follows
UTC ¼ CT þ z
where z is the number of standard time zones by which the observer is displaced to
the west of the Greenwich meridian. Some of the Earth time zones are indicated
below:
• Western European Time (0 h of difference with respect to UTC);
• Central European Time (+1 h);
• USA, comprising in turn:
– Atlantic Standard Time (−4 h);
– Eastern Standard Time (−5 h);
– Central Standard Time (−6 h);
– Mountain Standard Time (−7 h); and
– Pacific Standard Time (−8 h);
• Moscow Time (+3 h);
• Tokyo Time (+9 h).
In particular, there are four standard time zones in the conterminous USA. From
east to west they are: Eastern Standard Time (EST), Central Standard Time (CST),
Mountain Standard Time (MST) and Pacific Standard Time (PST), as shown in the
following figure (due to the courtesy of the National Atlas of the United States [59]).
Consequently, for any given place within the conterminous USA, UTC can be
computed by means of one of the following expressions
UTC ¼ EST þ 5
UTC ¼ CST þ 6
UTC ¼ MST þ 7
UTC ¼ PST þ 8
For example, we want to compute the UTC corresponding to
1:21:15 pm EST
First, we convert this time to a 24-h clock, as follows
1:21:15 pm EST þ 12h ¼ 13:21:15 EST
Then, adding 5 h (see above) to 13:21:15 EST, we have
13:21:15 EST þ 5h ¼ 18h 21m 15s UTC
The International System of Units (usually abbreviated as SI) defines the second
as the duration of 9,192,631,770 periods of the radiation corresponding to the
transition between the two hyperfine levels of the ground state of the Caesium 133
atom. This defines, in the abstract, the Atomic Time [38]. In practice, since atomic
clocks do not agree fully with one another, the weighted mean of several hundred
atomic clocks, placed in various laboratories on the Earth, is used to define the
so-called International Atomic Time (usually abbreviated as TAI).
As has been shown above, at about 1930, the Earth rotational period was found
to be irregular and therefore, for purposes of orbital calculations, time based on
Earth rotation was abandoned to choose a more uniform time scale based on the
Earth orbit about the Sun. Thus, the Ephemeris Time was defined as the time scale
which, together with the laws of motion, predicts correctly the positions of celestial
bodies. Because of this property, the Ephemeris Time is used as the argument of the
ephemerides, that is, of the tables giving the positions of the Sun, Moon, planets,
and respective satellites as a function of time. For this purpose, in 1958, the
International Astronomical Union (IAU) decided that “Ephemeris Time is reckoned
from the instant, near the beginning of the calendar year AD 1900, when the
geometric mean longitude of the Sun was 279 degrees 41 min 48.04 s, at which
instant the measure of Ephemeris Time was 1900 January 0, 12 o’clock precisely”.
Ephemeris Time was used for ephemeris calculations for the solar system until
1979, when it was replaced by Terrestrial Dynamical Time (TDT). TDT takes into
account relativistic effects and is based on International Atomic Time (TAI), which
has been defined above. To ensure continuity with Ephemeris Time, TDT was
defined to match ET for the date 1977 January 1. In 1991, the IAU modified the
definition of TDT to make it more precise. It was also renamed Terrestrial Time
(TT); however, the old name (Terrestrial Dynamical Time) is still used. According
to the United States Naval Observatory [73], Terrestrial Time is effectively equal to
International Atomic Time plus 32.184 s exactly. Thus, the epoch designated
“J2000.0” is specified as Julian date (see definition below) 2451545.0 TT, or 2000
January 1, 12 h TT. This epoch is also expressed as 2000 January 1, 11:59:27.816
TAI, or 2000 January 1, 11:58:55.816 UTC [74]. As a result of the decrease, at an
irregular rate, of the rotational velocity of the Earth, the difference
DT ¼ TT UT1
also decreases irregularly. The exact value of DT cannot be predicted and can only
result from the historical record and observations.
Since this value must be known to predict the correct times of astronomical
events (such as eclipses), then a series of polynomial expressions have been created
by Morrison and Stephenson [55], Espenak and Meeus [23], and Islam et al. [40] to
evaluate approximately DT during intervals of time of interest. Some examples
(from Ref. [23]) of such evaluations are given below. Let the decimal year (y) be
defined as follows
month 0:5
y ¼ year þ
12
which expression gives y for the middle of the month of interest.

Let t = y − 2000. Then, the approximate value of DT (in seconds) is given by:
DT ¼ 63:86 þ 0:3345 t 0:060374 t2 þ 0:0017275 t3 þ 0:000651814 t4

þ 0:00002373599t5 ðwith year between 1986 and 2005Þ
DT ¼ 62:92 þ 0:32217t þ 0:005589t2 ðwith year between 2005 and 2050Þ
DT ¼ 20 þ 32 ½ðy 1820Þ=1002 0:5629 ð2150 yÞ
ðwith year between 2050 and 2150Þ
Section 2.3 has shown that a change of reference system from the
topocentric-horizon (or UVW) system to the geocentric-equatorial (or XYZ) system
requires the latitude (u) and longitude (k) of the observer and the Greenwich
sidereal time (hG). If hG were known on a particular day and at a particular time,
then hG could be determined for any future day and time, because we know, as has
been shown above, that the Earth turns through 1.00273790935 rotations on its axis
per day.
Let hG0 be the value of hG relating to 0 h UT1 on the 1st of January of a
particular year. Let day 0 designate the 1st of January of the chosen year. Then, day
30 will designate the 31st of January, day 58 will designate the 28th of February, and
so on, as shown in the following table, where the numbers in parentheses refer to
the leap years.
Date Day No. Date Day No.

31 January 30 31 July 211(212)
28 February 58 31 August 242(243)
31 March 89(90) 30 September 272(273)
30 April 119(120) 31 October 303(304)
31 May 150(151) 30 November 333(334)
30 June 180(181) 31 December 364(365)
In addition, any time can be expressed as a decimal fraction of a day. By so

doing, any given set of values specifying a date and a time can be converted into a
single floating-point number (D), such that its integral part indicates the number of
days and its fractional part indicates the fraction of days elapsed from the chosen
origin of times.
Consequently, hG is expressible in degrees as follows
hG ¼ hG 0 þ 1:00273790935 360 D
or in radians as follows
hG ¼ hG0 þ 1:00273790935 2p D
The methods used to determine sidereal time are based on a continuous count of
days and fractions of day obtained by means of the Julian Day. The Julian Day or
Julian Day Number (JDN) is the number of days that have elapsed since an initial
epoch. The initial epoch has been set at noon Universal Time, Monday, the 1st of
January 4713 BC in the proleptic Julian Calendar (that is, in the Julian calendar
extended earlier in time so as to include dates preceding 4 AD). The initial epoch,
which is counted as Julian Day 0, corresponds to the 24th of November 4714 BC in
the proleptic Gregorian calendar (that is, in the calendar obtained by extending the
Gregorian calendar backwards to years earlier than 1582, using the Gregorian leap
year rules). The Julian date (JD) is a continuous count of days and fractions thereof
elapsed since the same initial epoch, as follows
UT1
JD ¼ JDN þ
24
According to the definition given above, a Julian date comprises an integral part
and a fractional part. The integral part of a Julian date is the Julian Day Number,
whereas its fractional part is the time of day elapsed from noon UT1 expressed as a
decimal fraction of a day.
Examples of this fractional part are given below. A fractional part equal to 0.5
indicates midnight UT1, because the Julian Day begins at noon. A fractional part of
0.1 indicates UT1/24 = 0.1, that is, UT1 = 2.4 h (or 2.4 60 = 144 min or
2.4 60 60 = 8640 s) elapsed from noon. Methods and tables for computing
JDN from a date expressed in year, month, and day of our Gregorian calendar have
been indicated by several authors and can also be found in the Internet. Curtis [20]
and Boulet [12] use the following method
JDN ¼ 367 y INTf1:75½y þ INTðm=12 þ 0:75Þg þ INTð275m=9Þ

þ d þ 1721013:5
where the function INT(x) means truncation (that is, retaining the integral part and
dropping the fractional part of the argument x of the function INT), y is the full
four-digit year, m is the month and d is the day of the Gregorian calendar date to be
converted. Curtis points out that y, m, and d are integers such that:
1901 y 2099
1 m 12
1 d 31
This means that the formula given above holds only if the year of interest falls
between 1901 and 2099.
These limitations do not affect another method, proposed by Jefferys [41], which
method is based on the following rules. As is the case with Boulet’s method,
Jefferys expresses a Gregorian calendar date as y-m-d, where y is the year, m is the
month number (January = 1, February = 2, etc.), and d is the day in the month. If
the month is January or February, then 1 must be subtracted from the year to get a
new value of y, and 12 must be added to the month to get a new value of m. Thus,
January and February of a given year are considered as being, respectively, the 13th
and the 14th month of the previous year.
Then, the computation follows the following scheme
a ¼ INT ðy=100Þ
b ¼ INTða=4Þ
c ¼ 2 aþb
e ¼ INT½365:25 ðy þ 4716Þ
f ¼ INT½30:6001 ðm þ 1Þ
JDN ¼ c þ d þ e þ f 1524:5
This is the Julian Day Number for the beginning of the desired date at 0 h, UT1.
To convert a Julian Day Number to a Gregorian calendar date, Jefferys [41] uses the
following method, assuming that the JDN is for 0 h, UT1 (so that it ends in 0.5).
The necessary calculations are shown below. Jefferys notes that his method does
not give dates accurately in the proleptic Gregorian Calendar; in particular, the
method fails if y is less than 400.
z ¼ JDN þ 0:5
w ¼ INT½ðz 1867216:25Þ=36524:25
x ¼ INT ðw=4Þ
a ¼ zþ1þw x
b ¼ a þ 1524
c ¼ INT½ðb 122:1Þ=365:25
d ¼ INT ð365:25cÞ
e ¼ INT½ðb dÞ=30:6001
f ¼ INT ð30:6001eÞ
day of month ¼ b d f
month ¼ e 1 or e 13ðthe number obtained must be less than or equal to 12Þ
year ¼ c 4715 ðif the month is January or FebruaryÞ or c 4716 ðotherwiseÞ
The U.S. Naval Observatory has an Internet-based calculator [71], which makes
it possible to compute automatically the Julian date.
As an example, let us convert the following Gregorian calendar date and time
3 February 2011 at 13:15:18 UT1
into the corresponding Julian date, by using Boulet’s method.

Since the condition 1901 y 2099 is satisfied for y = 2011, then we can
compute the Julian Day Number as follows
JDN ¼ 367 2011 INTf1:75 ½2011 þ INTð2=12 þ 0:75Þg þ INTð275 2=9Þ

þ 3 þ 1721013:5
¼ 738037 INTf1:75 ½2011 þ INTð0:917Þg þ INTð61:111Þ þ 3 þ 1721015:5
¼ 738037 INT½1:75 ð2011 þ 0Þ þ 61 þ 3 þ 1721013:5
¼ 738037 INTð3519:25Þ þ 1721077:5 ¼ 2455595:5 days
The universal time, expressed in hours, is
15 18
UT1 ¼ 13 þ þ ¼ 13:255
60 3600
Finally, the desired Julian date is
UT1 13:255
JD ¼ JDN þ ¼ 2455595:5 þ 2455596:05229 days
24 24
Now, let us apply Jefferys’ method to the same example. Since, in this case, the
month is February, then 1 must be subtracted from the year (2011) to get a new
value of y (2011 − 1 = 2010), and 12 must be added to the month (2) to get a new
value of m (2 + 12 = 14). Then the Julian Day Number is computed as follows
a ¼ INTðy=100Þ ¼ INTð2010=100Þ ¼ INTð20:10Þ ¼ 20

b ¼ INTða=4Þ ¼ INTð20=4Þ ¼ INTð20=4Þ ¼ 5
c ¼ 2 a þ b ¼ 2 20 þ 5 ¼ 13
e ¼ INT½365:25 ðy þ 4716Þ ¼ INT½365:25 ð2010 þ 4716Þ ¼ INTð2456671:5Þ
¼ 2456671
f ¼ INT½30:6001 ðm þ 1Þ ¼ INT½30:6001 ð14 þ 1Þ ¼ INTð459:0015Þ ¼ 459
JDN ¼ c þ d þ e þ f 1524:5 ¼ 13 þ 3 þ 2456671 þ 459 1524:5 ¼ 2455595:5
which is the same value as that computed previously by means of Boulet’s method.
Now, Jefferys’ method is applied to the case JDN = 2455595.5 to obtain the
corresponding Gregorian calendar date. Using the rules indicated above, we have
z ¼ 2455595:5 þ 0:5 ¼ 2455596

w ¼ INT½ð2455596 1867216:25Þ=36524:25 ¼ INTð16:109Þ ¼ 16
x ¼ INTð16=4Þ ¼ INTð4Þ ¼ 4
a ¼ 2455596 þ 1 þ 16 4 ¼ 2455609
b ¼ 2455609 þ 1524 ¼ 2457133
c ¼ INT½ð2457133 122:1Þ=365:25 ¼ INTð6726:929Þ ¼ 6726
d ¼ INTð365:25 6726Þ ¼ INTð2456671:5Þ ¼ 2456671
e ¼ INT½ð2457133 2456671Þ=30:6001 ¼ INTð15:098Þ ¼ 15
f ¼ INTð30:6001 15Þ ¼ INTð459:0015Þ ¼ 459
day of month ¼ 2457133 2456671 459 ¼ 3
month ¼ 15 13 ¼ 2
year ¼ 6726 4715 ¼ 2011
As was expected, the resulting Gregorian calendar date is 3 February 2011.

The U.S. Naval Observatory also gives algorithms in FORTRAN programming
language, which are due to Fliegel and van Flandern [25]. The current Julian epoch
has been set to the 1st of January 2000 at noon. This epoch is denoted by J2000.0
and corresponds to the Julian Day Number 2451545.0.
A Julian year has 365.25 days and consequently a Julian century has
36525 days.
Let T0 be the time, expressed in Julian centuries, elapsed between a given Julian
day J0 and J2000.0. Then the time T0 results from
J0 2451545:0
T0 ¼
36525
The time hG0 (that is, the Greenwich sidereal time at 0 h UT1) can be expressed
(in seconds) as a function of the dimensionless time T0 by means of the following
formula given by Aoki et al. [4]:
hG0 ¼ 24110s :54841 þ 8640184s :812866 T0 þ 0s :093104T0 2 0s :000006210T0 3
where the superscript s stands for seconds. If, as is often the case, we want to
express hG0 in degrees, then the coefficients appearing in the formula given above
must be multiplied by 360/(24 3600). The same formula expressed in degrees is
hG0 ¼ 100:460618375 þ 36000:7700536 T0 þ 0:000387933 T0 2 2:5875

108 T0 3
Since the value computed by means of the preceding formula may be outside the
interval 0° < hG0 < 360°, then the computed value must, if necessary, be brought
into that interval by adding or subtracting an integral multiple of 360°.
This done, the Greenwich sidereal time (hG) relating to any other universal time
than 0 h can be found as follows
UT1
hG ¼ hG0 þ 360:985647366
24
where 360.985647366 = 1.00273790935 360 is the number of degrees covered

by the Earth in its rotation about its axis in 24 h (solar time). If the (local) sidereal
time h is required for a site placed at an east longitude k, then
h ¼ hG þ k
As shown above, in case of the value of h being greater than 360°, it is necessary
to bring it into the interval 0° < h < 360° by subtracting an integral number of 360°
from the computed value.
As an example of application, let us compute the local sidereal time in degrees
for Kiruna, Sweden (latitude u = 65°.85N; longitude k = 20°.2167E) on the 13th of
February 2012 at 2:30:00 UT1. First, we use Jefferys’ method to compute the Julian
Day Number, that is, the Julian date relating to 13th of February 2012 at 00:00:00
UT1. Since, in this case, the month is February, then 1 must be subtracted from the
year (2012) to get a new value of y (2012 − 1 = 2011), and 12 must be added to the
month (2) to get a new value of m (2 + 12 = 14).
Then the Julian Day Number (J0) is computed as follows

b ¼ INTða=4Þ ¼ INTð20=4Þ ¼ INTð5Þ ¼ 5
c ¼ 2 a þ b ¼ 2 20 þ 5 ¼ 13
e ¼ INT½365:25 ðy þ 4716Þ ¼ INT½365:25 ð2011 þ 4716Þ ¼ INTð2457036:75Þ
¼ 2457036
f ¼ INT½30:6001 ðm þ 1Þ ¼ INT½30:6001 ð14 þ 1Þ ¼ INTð459:0015Þ ¼ 459
J0 ¼ c þ d þ e þ f 1524:5 ¼ 13 þ 13 þ 2457036 þ 459 1524:5 ¼ 2455970:5
Second, the dimensionless time T0 results from the expression shown above
J0 2451545:0 2455970:5 2451545:0

T0 ¼ ¼ ¼ 0:12116358658453
36525 36525
Thirdly, hG0 is computed by using Aoki’s formula expressed in degrees
hG 0 ¼ 100:4606184 þ 36000:77005 T0 þ 0:000387933 T0 2 2:5875 108 T0 3
Substituting T0 = 0.12116358658453 in the preceding formula, there results
hG0 ¼ 100:4606184 þ 36000:77005 0:12116358658453 þ 0:000387933

0:121163586584532 2:5875 108 0:121163586584533
¼ 4462 :44304315801945
Since this value is outside the interval [0, 360], we bring it into that interval by
subtracting a multiple of 360° from it. To this end, we observe that

4462:44304315801945
INT ¼ 12
360
It follows that
hG0 ¼ 4462:44304315801945 12 360 ¼ 142 :44304315801945
The universal time given in this example is 2:30:00, that is,
UT1 ¼ 2:5 h
Thus, the Greenwich sidereal time hG is computed by replacing hG0 with the
value obtained above and UT1 with 2.5 into the expression
UT1
hG ¼ hG0 þ 360:985647366
24
This yields
hG ¼ 142 :44304315801945 þ 360:985647366 2:5=24

¼ 180 :04571475864445
Finally, the east longitude of Kiruna (k = 20°.2167) must be added to hG to

obtain the local sidereal time, as follows
h ¼ hG þ k ¼ 180 :04571475864445 þ 20 :2167 ¼ 200 :26241475864445
2.5 Orbital Elements from Angle and Range

Measurements
The orbital elements of a space object Q revolving around the Earth and shown in
the following figure are determined when its position (r) and velocity (r′) vectors,
with respect to the geocentric-equatorial system XYZ, are known at a given time.
Section 2.3 has shown how to determine Q as a function of the line-of-sight

vector q of Q with respect to the topocentric-horizon system UVW (located at the
radar station P) and the position vector rE of P with respect to the
geocentric-equatorial system. The geocentric position vector r OQ of the space
object Q results from
r ¼ rE þ q ¼ rE þ quq
where uq q/q is the unit vector having the direction of q. The velocity (r′) and
acceleration (r″) vectors of Q, with respect to the geocentric-equatorial system XYZ,
result from differentiating once and twice the preceding expression with respect to
time. This yields
r0 ¼ r0E þ q0 uq þ qu0q
r00 ¼ r00E þ q00 uq þ q0 u0q þ q0 u0q þ qu00q ¼ r00E þ q00 uq þ 2q0 u0q þ qu00q
Now, all vectors appearing in the preceding expressions must be expressed in the
same reference system, that is, in the geocentric-equatorial reference system XYZ.
To this end, it is to be noted that rE OP is a vector which rotates with the Earth at
a constant angular velocity xE equal to
xE ¼ xE uZ
where xE is the magnitude of xE, and uZ is the unit vector of the Z-axis. Since rE
rotates with the Earth, then its first (r′E) and second (r″E) time derivatives result
from the following expressions
r0E ¼ xE rE
r00E ¼ xE ðxE rE Þ
Let ‘X, ‘Y, and ‘Z be the direction cosines of PQ with respect to the
topocentric-equatorial system XtYtZt, having its origin in P and its axes parallel to,
respectively, X, Y, and Z. The unit vector uq = q/q is expressible as follows
uq ¼ ‘X uX þ ‘Y uY þ ‘Z uZ
Since the unit vectors uX, uY, and uZ do not change with time, then the first and
second time derivatives of the preceding expression are
u0q ¼ ‘0X uX þ ‘0Y uY þ ‘0Z uZ

u00q ¼ ‘00X uX þ ‘00Y uY þ ‘00Z uZ
Now, since uq = (cos d cos a) uX + (cos d sin a) uY + (sin d) uZ, where d and a
are, respectively, the declination and right ascension of Q, then the direction cosines
of PQ with respect to the topocentric-equatorial system are expressible in terms of d
and a as follows
2.5 Orbital Elements from Angle and Range Measurements 165
‘X ¼ cos d cos a
‘Y ¼ cos d sin a
‘Z ¼ sin d
Differentiating the preceding expressions with respect to time yields
‘0X ¼ a0 sin a cos d d0 cos a sin d

‘0Y ¼ a0 cos a cos d d0 sin a sin d
‘0Z ¼ d0 cos d
‘00X ¼ a00 sin a cos d d00 cos a sin d ða0 2 þ d0 2 Þ cos a cos d þ 2a0 d0 sin a sin d
‘00Y ¼ a00 cos a cos d d00 sin a sin d ða0 2 þ d0 2 Þ sin a cos d 2a0 d0 cos a sin d
‘00Z ¼ d00 cos d d0 2 sin d
Now, let ‘U, ‘V, and ‘W be the direction cosines of PQ with respect to the
topocentric-horizon system UVW, having its origin in P and its axes pointing to,
respectively, east, north, and zenith.
As shown in Sect. 2.3, the unit vector of the direction PQ is
uq ¼ ðcos h sin AÞuU þ ðcos h cos AÞuV þ ðsin hÞ uW
where h and A are, respectively, the altitude angle and the azimuth angle of Q.
Consequently, the direction cosines of PQ with respect to the
topocentric-horizon system UVW are
‘U ¼ cos h sin A
‘V ¼ cos h cos A
‘W ¼ sin h
The direction cosines ‘X, ‘Y, and ‘Z of PQ with respect to the

geocentric-equatorial system XYZ can be obtained by means of the co-ordinate
transformation shown in Sect. 2.3, that is,
2 3 2 32 3
‘X sin h sin u cos h cos u cos h ‘U
4 ‘Y 5 ¼ 4 cos h sin u sin h cos u sin h 54 ‘V 5:
‘Z 0 cos u sin u ‘W
replacing ‘X, ‘Y, ‘Z, ‘U, ‘V, and ‘W by their respective values yields
2 3 2 32 3
cos d cos a sin h sin u cos h cos u cos h cos h sin A
4 cos d sin a 5 ¼ 4 cos h sin u sin h cos u sin h 54 cos h cos A 5
sin d 0 cos u sin u sin h
Expanding the matrix product in the preceding equality leads to
cos d cos a ¼ sin h cos h sin A sin u cos h cos h cos A þ cos u cos h sin h
hence
ðcos u sin h sin u cos h cos AÞ cos h sin h cos h sin A

cos a ¼
cos d
cos d sin a ¼ cos h cos h sin A sin u sin h cos h cos A þ cos u sin h sin h
hence
ðcos u sin h sin u cos h cos AÞ sin h þ cos h cos h sin A

sin a ¼
cos d
sin d ¼ cos u cos h cos A þ sin u sin h
The preceding expressions of cos a, sin a, and sin d can be simplified by using
the hour angle H, defined in Sect. 2.4, which is the angular distance between the
object observed and the local meridian. In terms of the variables used above, the
hour angle is expressible as follows
H ¼ha
The sine and cosine of the hour angle can be expressed as follows
sinðh aÞ ¼ sin h cos a cos h sin a

cosðh aÞ ¼ cos h cos a þ sin h sin a
Thus, substituting the following expressions

cos a ¼
cos d
sin a ¼
cos d
into
sinðh aÞ ¼ sin h cos a cos h sin a
yields

sinðh aÞ ¼ sin h
cos d

cos h
cos d
cos h sin A
¼
cos d
Likewise, substituting

cos a ¼
cos d
sin a ¼
cos d
into
cosðh aÞ ¼ cos h cos a þ sin h sin a
yields

cosðh aÞ ¼ cos h
cos d

þ sin h
cos d
cos u sin h sin u cos h cos A
¼
cos d
This expression makes it possible to compute the hour angle, as follows

H ¼ arccos
cos d
which holds if sin H > 0. Otherwise, if sin H < 0, then

H ¼ 2p arccos
cos d
Since the altitude (h) and declination (d) angles range from −p/2 to p/2 radians,
then neither cos h nor cos d can be negative. Consequently, the expression
cos h sin A
sin H ¼
cos d
shows that sin H is negative when sin A is positive, which happens when A ranges
from 0 to p radians. In summary, let the topocentric angles azimuth (A) and altitude
(h) of a space object be known at a given time. Let also the latitude (u) of the
tracking station and the sidereal time (h) be known at the same time. Then, the
topocentric declination (d) results from the equation derived above
which, solved for d, yields
d ¼ arcsin ðcos u cos h cos A þ sin u sin hÞ
This done, the hour angle (H) results from

2p arccos ½ðcos u sin h sin u cos h cos AÞ= cos d ð0\A\pÞ
H¼
arccos½ðcos u sin h sin u cos h cos AÞ= cos d ðp A 2pÞ
and the right ascension (a) results from
a¼hH
If the azimuth and altitude angles are known as functions of time, then the right
ascension and declination angles can be computed as functions of time by means of
the expressions given above. Then, these functions are differentiated with respect to
time, and the results are introduced into the following expressions
u00q ¼ ‘00X uX þ ‘00Y uY þ ‘00Z uZ
where
‘X ¼ cos d cos a
‘Y ¼ cos d sin a
‘Z ¼ sin d

‘0Z ¼ d0 cos d
‘00X ¼ a00 sin a cos d d00 cos a sin d ða0 2 þ d0 2 Þ cos a cos d þ 2a0 d0 sin a sin d
‘00Y ¼ a00 cos a cos d d00 sin a sin d ða0 2 þ d0 2 Þ sin a cos d 2a0 d0 cos a sin d
‘00Z ¼ d00 cos d d0 2 sin d
This makes it possible to compute the unit vector uq and its time derivatives u′q
and u″q. In order to compute a′ and d′ from A′ and h′, the expression
is differentiated with respect to time, taking into account that u = constant.

This yields
d0 cos d ¼ h0 cos u sin h cos A A0 cos u cos h sin A þ h0 sin u cos h
which in turn, solved for d′, yields
A0 cos u cos h sin A þ h0 ðsin u cos h cos u sin h cos AÞ

d0 ¼
cos d
Now, the expression sin H = −(cos h sin A)/cos d is differentiated with respect to
time. This yields
h0 sin h sin A A0 cos h cos A d0 cos h sin A sin d

H 0 cos H ¼
cos d cos d cos2 d
ðA cos h cos A h sin h sin AÞ cos d þ d0 cos h sin A sin d
0 0
¼
cos2 d
Since

cos H ¼
cos d
then substituting this expression of cos H into the expression of H′ cos H yields

0 cos u sin h sin u cos h cos A
0
H cos H ¼ H
cos d
ðA0 cos h cos A h0 sin h sin AÞ cos d þ d0 cos h sin A sin d
¼
cos2 d
The preceding equation, solved for H′, yields
A0 cos h cos A h0 sin h sin A þ d0 cos h sin A tan d

H0 ¼
Now, since H = h − a, then
H 0 ¼ h0 a0 ¼ x E a0
where xE is the angular velocity of the Earth about its axis. It follows that
a0 ¼ x E H 0
that is,
A0 cos h cos A h0 sin h sin A þ d0 cos h sin A tan d

a0 ¼ xE þ
As an example of application, we compute the classical orbital elements of a

space object, which the Catalina station of the University of Arizona, Tucson
(latitude u = 32°.417 N, longitude k = −110°.732 E, and height H = 2509 m on
the mean sea level) detected on the 5th of April 2004 at 6:00:00 UT1, obtaining the
following data:
Slant range q ¼ 27148 km

Azimuth A ¼ 128
Altitude h ¼ 41
Range rate q0 ¼ 2:267 km=s
Azimuth rate A0 ¼ 1:86 105 rad=s
Altitude rate h0 ¼ 3:42 105 rad=s
The three components of the position vector (rE) of the tracking station, with
respect to the geocentric-equatorial system XYZ, are
where uX, uY, and uZ are the unit vectors along, respectively, X, Y, and Z, xH and zH
are the two co-ordinates of the point PH representing the position of the tracking
station. Let us compute the local sidereal time h by means of the date
5 April 2004 at 6:00:00 UT1

The Julian Day Number (J0) is computed as follows

b ¼ INTða=4Þ ¼ INT ð20=4Þ ¼ INT ð5Þ ¼ 5
c ¼ 2 a þ b ¼ 2 20 þ 5 ¼ 13
e ¼ INT½365:25 ðy þ 4716Þ ¼ INT½365:25 ð2004 þ 4716Þ ¼ 2454480
f ¼ INT½30:6001 ðm þ 1Þ ¼ INT½30:6001 ð4 þ 1Þ ¼ INT ð153:0005Þ ¼ 153
J0 ¼ c þ d þ e þ f 1524:5 ¼ 13 þ 5 þ 2454480 þ 153 1524:5 ¼ 2453100:5
The dimensionless time T0 results from the following expression
J0 2451545:0 2453100:5 2451545:0

T0 ¼ ¼ ¼ 0:042587
36525 36525
Now, hG0 is computed by using Aoki’s formula
hG0 ¼ 100:4606184 þ 36000:77005 T0 þ 0:000387933 T0 2 2:5875 108 T0 3

¼ 100:4606184 þ 36000:77005 0:042587 þ 0:000387933 0:0425872
2:5875 108 0:0425873 ¼ 1633 :625413
Since this value is outside the interval [0, 360], we bring it into that interval by
subtracting a multiple of 360° from it. To this end, we observe that

1633:625
INT ¼4
360
It follows that
hG0 ¼ 1633:625413 4 360 ¼ 193 :6254128
The universal time given in this example is 6:00:00, that is, UT1 = 6.0 h.
Thus, the Greenwich sidereal time hG is computed by replacing hG0 with
193.6254128 and UT1 with 6.0 into the following expression

UT1
hG ¼ hG0 þ 360:985647366
24
This yields

6
hG ¼ 193:6254128 þ 360:985647366 ¼ 283 :8718246
24
The east longitude of Catalina (k = −110°.732 east) must be added to hG to

obtain the local sidereal time, as follows
h ¼ hG þ k ¼ 283 :8718246 110 :732 173 :14
Now, xH and zH are computed as follows

8 9
< aE =
xH ¼ þ cos u
: 1 ð2f f 2 Þ sin2 u 12
H
;
8 9
< aE ð 1 f Þ 2 =
zH ¼ þ sin u
: 1 ð2f f 2 Þ sin2 u 12
H
;
where aE = 6378.137 km and f = 0.0033528 are, respectively, the equatorial radius

and the flattening of the Earth (represented as an ellipsoid), and u is the geodetic
latitude (that is, the angle between the equatorial plane and the local vertical)
measured at the tracking station. Taking u = 32°.417, xH and zH result from
n 1=2 o
xH ¼ 6378:137= 1 2 0:0033528 0:00335282 sin2 32 :417 þ 2:509
cos 32 :417 ¼ 5391:552 km
n
zH ¼ 6378:137 ð1 0:0033528Þ2 = 1 2 0:0033528 0:00335282
1=2 o
sin2 32 :417 þ 2:509 sin 32 :417 ¼ 3400:921 km
Hence, the position vector rE of the tracking station, with respect to the
geocentric-equatorial system XYZ, is

¼ ð5391:552 cos 173 :14ÞuX þ ð5391:552 sin 173 :14Þ uY þ 3400:921uZ
¼ 5352:954 uX þ 643:987 uY þ 3400:921 uZ ðkmÞ
Let us compute now the declination d of the observed object with respect to the
topocentric-equatorial system XtYtZt, by means of the following expression
d ¼ arcsin ðcos u cos h cos A þ sin u sin hÞ

¼ arcsinðcos 32 :417 cos 41 cos 128 þ sin 32 :417 sin 41 Þ ¼ 2 :323
The given azimuth (A = 128°) lies between 0° and 180°. Consequently, the hour
angle H results from the following expression

H ¼ 360 arccos
cos d

cos 32 :417 sin 41 sin 32 :417 cos 41 cos 128

¼ 360 arccos
cosð2 :323Þ
¼ 323 :472
Thus, the right ascension of the observed object results from
a ¼ h H ¼ 173 :14 323 :472 ¼ 150 :332
Now, we compute the unit vector uq of the line joining the tracking station with
the object, with respect to the topocentric-equatorial system XtYtZt, by means of
where
‘X ¼ cos d cos a
‘Y ¼ cos d sin a
‘Z ¼ sin d
This yields
uq ¼ ½cosð2 :323Þ cosð150 :332ÞuX þ ½cosð2 :323Þ sinð150 :332ÞuY

þ ½sinð2 :323ÞuZ ¼ 0:8682 uX 0:4946 uY 0:0405 uZ
The position vector r (in km) of the observed object results from
r ¼ rE þ quq ¼ 5352:954 uX þ 643:987 uY þ 3400:921 uZ

þ 27148 ð0:8682 uX 0:4946 uY 0:0405 uZ Þ
¼ 28922:848 uX 12783:414 uY þ 2301:427 uZ
The velocity vector r′E of the tracking station, with respect to the celestial
geocentric-equatorial system XYZ, results from
r0E ¼ xE rE
where rE = −5352.954 uX + 643.987 uY + 3400.921 uZ is the position vector of the

tracking station and xE = xE uZ is the angular velocity of the Earth around its axis
with respect to XYZ. Remembering that xE = 7.292 10−5 rad/s, there results
2 3
uX uY uZ
6 7
r0E ¼ xE rE ¼ 4 0 0 7:292 105 5
5352:954
643:987 3400:921

¼ 643:987 7:292 10 uX þ ð5352:954 7:292 105 ÞuY
5
¼ 0:047 uX 0:390 uY ðkm=sÞ
Now, the declination rate d′ of the observed object results from
d0 ¼ ½A0 cos u cos h sin A þ h0 ðsin u cos h cos u sin h cos AÞ= cos d

¼ 1:86 105 cos 32 :417 cos 41 sin 128 þ 3:42 105 ðsin 32 :417 cos 41
cos 32 :417 sin 41 cos 128 Þ= cosð2 :323Þ ¼ 3:486 105 ðrad=sÞ
The right ascension rate a′ of the observed object results from
a0 ¼ xE þ ðA0 cos h cos A h0 sin h sin A þ d0 cos h sin A tan dÞ=ðcos u sin h
sin u cos h cos AÞ ¼ 7:292 105 þ ½1:86 105 cos 41 cos 128
3:42 105 sin 41 sin 128 þ 3:486 105 cos 41 sin 128 tan ð2 :323Þ
=½cos 32 :417 sin 41 sin 32 :417 cos 41 cos 128 ¼ 6:062 105 ðrad=sÞ
The direction cosine rate vector u′q results from
where

‘0Z ¼ d0 cos d
Hence,
‘0X ¼ 6:062 105 sinð150 :332Þ cosð2 :323Þ 3:486 105

cosð150 :332Þ sinð2 :323Þ ¼ 2:875 105
‘0Y ¼ 6:062 105 cosð150 :332Þ cosð2 :323Þ 3:486 105
sinð150 :332Þ sinð2 :323Þ ¼ 5:333 105
‘0Z ¼ 3:486 105 cosð2 :323Þ ¼ 3:483 105
which, substituted into u′q = ‘′X uX + ‘′Y uY + ‘′Z uZ, yield
u0q ¼ 2:875 105 uX 5:333 105 uY þ 3:483 105 uZ ðrad=sÞ
The velocity vector r′ of the observed object, with respect to the

geocentric-equatorial reference system XYZ, results from
r0 ¼ r0E þ q0 uq þ qu0q
Hence, the velocity vector r′ (in km/s) of the observed object is
r0 ¼ 0:047 uX 0:390 uY þ 2:267 ð0:8682 uX 0:4946 uY 0:0405 uZ Þ

þ 27148 ð2:875 105 uX 5:333 105 uY þ 3:483 105 uZ Þ
¼ 1:235 uX 2:959 uY þ 0:854 uZ
In summary, the position (r) and velocity (r′) vectors of the space object at the
epoch of observation, with respect to the geocentric-equatorial system XYZ, are
r ¼ 28922:848 uX 12783:414 uY þ 2301:427 uZ ðkmÞ

r0 ¼ 1:235 uX 2:959 uY þ 0:854 uZ ðkm=sÞ
Now that the position and velocity vectors are known, the corresponding orbital
elements can be computed, as will be shown below.
The radius vector (in km) and the square of the velocity (in km2/s2) at epoch are
h i12
r0 ¼ ð28922:848Þ2 þ ð12783:414Þ2 þ 2301:4272 ¼ 31705:573
v20 ¼ ð1:235Þ2 þ ð2:959Þ2 þ 0:8542 ¼ 11:01
The vis-viva integral
v2 2 1
¼
l r a
makes it possible to compute the major semi-axis a (in km) as follows
1 1
a¼ ¼ ¼ 28201:744
2 v20 2 11:01

r0 l 31705:573 398600:4
The angular momentum vector per unit mass (in km2/s) is

2 3
uX uY uZ
6 7
h ¼ r r0 ¼ 4 28922:848 12783:414 2301:427 5
1:235 2:959 0:854
¼ 4107:113 uX þ 21857:85 uY þ 69795:191 uZ
and its magnitude (in km2/s) is
1
h i12
h ¼ ðh hÞ2 ¼ ð4107:113Þ2 þ 21857:852 þ 69795:1912 ¼ 73253:005
The semi-latus rectum p (in km) results from
h2 73253:0052
p¼ ¼ ¼ 13462:084
l 398600:4
The eccentricity e results from a = p(1 − e2), which, solved for e, yields
1
p12 13462:084 2
e¼ 1 ¼ 1 ¼ 0:72294582
a 28201:744
The inclination angle i (in degrees) of the orbit with respect to the equatorial
plane is

hZ 69795:191
i ¼ arccos ¼ arccos ¼ 17 :674578
h 73253:005
The eccentricity vector e results from
r0 h r
e¼
l r0
ðr0 hÞ=l ¼ ½ð2:959 69795:191 21857:85 0:854ÞuX þ ð1:235 69795:191

4107:113 0:854Þ uY þ ð1:235 21857:85 4107:113 2:959ÞuZ
=398600:4 ¼ 0:5649521uX þ 0:2074494uY 0:0982119 uZ
r=r0 ¼ ð28922:848 uX 12783:414 uY þ 2301:427 uZ Þ=31705:573

¼ 0:9122323 uX 0:4031914 uY þ 0:0725875 uZ
e ¼ ðr0 hÞ=l r=r0 ¼ ð0:5649521 þ 0:9122323Þ uX

þ ð0:2074494 þ 0:4031914Þ uY þ ð0:0982119 0:0725875Þ uZ
¼ 0:3427113 uX þ 0:6106408 uY 0:1707994 uZ
The node vector n is defined by n uZ h. In the present case, there results
n ¼ nX uX þ nY uY ¼ hY uX þ hX uY ¼ 21857:85 uX 4107:113 uY
The magnitude n of the node vector is

1 1
n ¼ ðn nÞ2 ¼ 21857:852 þ 4107:1132 2 ¼ 22240:368
The right ascension X of the ascending node (in degrees) is

n
X
X ¼ arccos ðnY 0Þ
n
n
X ¼ 360 arccos
X
ðnY \0Þ
n
In the present case, nY = −4107.113 < 0; thus,

21857:85
X ¼ 360 arccos ¼ 190 :64185
22240:368
The argument of perigee x results from

n e
x ¼ arccos ðeZ 0Þ
ne
n e
x ¼ 360 arccos ðeZ \0Þ
ne
In the present case, eZ = −0.1707994 < 0; hence,

21857:85 0:3427113 4107:113 0:6106408
x ¼ 360 arccos
22240:368 0:72294582
¼ 231 :54665
The true anomaly at epoch /0 results from

e r
/0 ¼ arccos ðr r0 0Þ
er0

e r
/0 ¼ 360 arccos ðr r0 \0Þ
er0
r r0 ¼ 28922:848 ð1:235Þ 12783:414 ð2:959Þ þ 2301:427 0:854

¼ 75511:258 [ 0
e r ¼ 0:3427113 ð28922:848Þ þ 0:6106408 ð12783:414Þ 0:1707994
2301:427 ¼ 18111:34334
er0 ¼ 0:72294582 31705:573 ¼ 22921:41147
Thus, the true anomaly (in degrees) at epoch is

18111:34334
/0 ¼ arccos ¼ 142 :19949
22921:41147
In summary, the object detected by the given radar station at the given time
revolves about the Earth in an elliptic orbit having the following elements
a ¼ 28201:744 km X ¼ 190 :64185

e ¼ 0:72294582 x ¼ 231 :54665
i ¼ 17 :674578 /0 ¼ 142 :19949

(Method of Gauss)
As has been shown in the preceding paragraphs, a set of six independent quantities
is required to determine the motion of a celestial body. These six quantities may be
the three components of the position vector and the three components of the
velocity vector, or the six classical elements. A radar station like that described in
Sect. 2.5 provides range and range-rate measurements. Thus, three range mea-
surements (range, declination, and azimuth) plus three range-rate measurements
(range rate, declination rate, and azimuth rate) provide the six required quantities.
This implies the availability of a Doppler radar.
By contrast, the present paragraph and the following one will show how to
determine the orbit of a celestial body when only angular observations are possible.
This happens when only a telescope can be used as a means of observation. When
only angular measurements are possible (e.g. the declination and the right ascension
of the body observed), then three distinct observations are required, each of which
provides the declination and the right ascension of the body.
Following Curtis [20] and Boulet [12], let t1, t2, and t3 be the three distinct times
at which the three single angular observations are performed. Let P1, P2, and P3 be
the three positions of the observed body at, respectively, t1, t2, and t3, as shown in
the following figure. Let r1, r2, and r3 be the three position vectors of the observed
2.6 Orbital Elements from Three Measurements of Angles (Method of Gauss) 179
body with respect to the geocentric-equatorial reference system XYZ. Let rE1, rE2,
and rE3 be the three position vectors of the point of observation with respect to the
same system of reference.
The relation between r1, r2, and r3 and rE1, rE2, and rE3 is
r1 ¼ rE1 þ q1 ¼ rE1 þ q1 ðq1 =q1 Þ ¼ rE1 þ q1 u1

where q1, q2, and q3 are the three position vectors of the observed body, with
respect to the topocentric-equatorial reference system XtYtZt, at, respectively, t1, t2,
and t3. Likewise, q1, q2, and q3 are the magnitudes of the vectors q1, q2, and q3, and
u1 = q1/q1, u2 = q2/q2, and u3 = q3/q3 are the three unit vectors along, respec-
tively, q1, q2, and q3. The three unit vectors u1, u2, and u3 are determined by
measuring the declination d and the right ascension a of the observed body at,
respectively, t1, t2, and t3, as follows
u1 ¼ ðcos d1 cos a1 Þ uX þ ðcos d1 sin a1 Þ uY þ ðsin d1 Þ uZ

where uX, uY, and uZ are the three unit vectors along, respectively, X, Y, and Z.
The three vector equations written above

provide nine scalar equations in twelve unknowns, which are the three components
of each of the three position vectors r1, r2, and r3 plus the three magnitudes q1, q2,
and q3 (3 3 + 3 = 12 unknowns). Three additional scalar equations are provided

by the fact that the motion of the observed body is confined in a plane, because the
moment of momentum per unit mass h is a constant vector (see Sect. 1.1).
Consequently, the three position vectors r1, r2, and r3 are coplanar. This means that
one of these vectors results from a linear combination of the other two. Without loss
of generality, we suppose r2 to be a linear combination of r1 and r3, so that
r 2 ¼ c1 r 1 þ c3 r 3
This equation, added to the 9 + 3 = 12 scalar equations written above, intro-

duces two new unknowns (c1 and c3). Thus, we have 12 scalar equations and
12 + 2 = 14 unknowns. In addition to the constancy of h, the Keplerian unper-
turbed motion of the observed body implies that the position vector of that body at
any time can be expressed in terms of the position (r) and velocity (v) vectors at any
other time by means of the Lagrangian coefficients (f and g). Thus, in the present
case, the position vectors r1 and r3 of the observed body at times, respectively, t1
and t3 can be expressed in terms of r2 and v2 at time t2, as follows
r1 ¼ f1 r2 þ g1 v2
r3 ¼ f3 r2 þ g3 v2
where f1 and g1 are the Lagrangian coefficients computed at time t1; likewise, f3 and
g3 are the Lagrangian coefficients computed at time t3.
As shown in Sect. 1.12, in case of small intervals of time between two con-
secutive observations, the Lagrangian coefficients f and g depend only on the
distance existing at the initial time between the attracted body and its centre of
attraction. In this case, designating the intermediate time t2 as the initial time and r2
as the distance existing at time t2 between the two mutually attracting bodies, the
Lagrangian coefficients f1, g1, f3, and g3 depend only on the distance r2. The two
vector equations
r1 ¼ f1 r2 þ g1 v2
r3 ¼ f3 r2 þ g3 v2
correspond to six scalar equations. We have then 12 + 6 = 18 equations.

On the other hand, the new unknowns introduced by these six equations are four:
the three components of the velocity vector v2 and the distance r2. We have then
14 + 4 = 18 unknowns. Thus, under the hypothesis made above, there is only one
solution to the problem of determining the vectors r2 and v2 at the initial time t2.
To this end, we first solve the equation
r 2 ¼ c1 r 1 þ c3 r 3
for c1 and c3. The vector product of all terms of this equation by r3 yields
r2 r3 ¼ c1 ðr1 r3 Þ þ c3 ðr3 r3 Þ
Since the vector product of any vector by itself is the zero vector, then c3(r3 r3)
vanishes identically. Thus,
r2 r3 ¼ c1 ðr1 r3 Þ
The scalar product of (r1 r3) by all terms of the preceding expression yields
ðr1 r3 Þ ðr2 r3 Þ ¼ c1 ðr1 r3 Þ ðr1 r3 Þ ¼ c1 jr1 r3 j2
where jr1 r3 j2 is the squared magnitude of r1 r3. The preceding equation,

solved for c1, yields
ðr1 r3 Þ ðr2 r3 Þ
c1 ¼
j r1 r3 j 2
By operating likewise, we form the vector product of all terms of the equation
r 2 ¼ c1 r 1 þ c3 r 3
by r1. This yields
r2 r1 ¼ c1 ðr1 r1 Þ þ c3 ðr3 r1 Þ ¼ c3 ðr3 r1 Þ
The scalar product of (r3 r1) and all terms of the preceding expression yields
ðr3 r1 Þ ðr2 r1 Þ ¼ c3 ðr3 r1 Þ ðr3 r1 Þ ¼ c3 jr3 r1 j2
Hence,
c3 ¼
j r3 r1 j 2
Now, we form the vector product r1 r3 and introduce
r1 ¼ f1 r2 þ g1 v2
r3 ¼ f3 r2 þ g3 v2
into the product r1 r3 . This yields
r1 r3 ¼ ðf1 r2 þ g1 v2 Þ ðf3 r2 þ g3 v2 Þ
¼ f1 f3 ðr2 r2 Þ þ f1 g3 ðr2 v2 Þ þ g1 f3 ðv2 r2 Þ þ g1 g3 ðv2 v2 Þ
Now, since the vector product of any vector by itself yields the zero vector, and
r 2 v2 ¼ h
v2 r2 ¼ ðr2 v2 Þ ¼ h
where h is the constant moment of momentum, then
r1 r3 ¼ ðf1 g3 g1 f3 Þh
r3 r1 ¼ ðr1 r3 Þ ¼ ðf1 g3 g1 f3 Þh
It follows that
jr1 r3 j2 ¼ jr3 r1 j2 ¼ ðf1 g3 g1 f3 Þ2 h2
Likewise, we form the vector product r2 r3 and introduce r3 = f3 r2 + g3 v2

into the product r2 r3. This yields
r2 r3 ¼ r2 ðf3 r2 þ g3 v2 Þ ¼ f3 ðr2 r2 Þ þ g3 ðr2 v2 Þ ¼ g3 h
Again, we form the vector product r2 r1 and introduce r1 = f1 r2 + g1 v2 into

the product r2 r1. This yields
r2 r1 ¼ r2 ðf1 r2 þ g1 v2 Þ ¼ f1 ðr2 r2 Þ þ g1 ðr2 v2 Þ ¼ g1 h
In summary, we have obtained the following expressions
r3 r1 ¼ ðr1 r3 Þ ¼ ðf1 g3 g1 f3 Þh
r 2 r 1 ¼ g1 h
which in turn, substituted into
c3 ¼
j r3 r1 j 2
yield
½ðf1 g3 g1 f3 Þh ½g1 h

c3 ¼
ðf1 g3 g1 f3 Þ2 h2
Since h h = h2, then

g1
c3 ¼
f 1 g3 g1 f 3
Let us consider now the expression
c1 ¼
j r1 r3 j 2
which has been derived above. Substituting
ðr1 r3 Þ ¼ ðf1 g3 g1 f3 Þh
r 2 r 3 ¼ g3 h
into this expression yields
½ðf1 g3 g1 f3 Þh ½g3 h

c1 ¼
ðf1 g3 g1 f3 Þ2 h2
that is,
g3
c1 ¼
f1 g3 g1 f3
By so doing, the coefficients c1 and c3 appearing in the equation
r 2 ¼ c1 r 1 þ c3 r 3
depend on the Lagrange coefficients f and g only. We set
s1 ¼ t 1 t 2
s3 ¼ t 3 t 2
where s1 and s3 are, by hypothesis, small intervals. Therefore, we can take the first
two terms of the series expansions for the Lagrangian coefficients f and g
1 1 1
f ¼ 1 0 s2 þ 0 k0 s3 þ 20 15k20 þ 3w0 s4 þ
2 2 24
1 3 1
g ¼ s 0 s þ 0 k0 s þ
4
6 4
(see Sect. 1.12). In other words, we truncate these series expansions after the first
two terms, as follows

1 l 2 1 l 2
f1 1 s f3 1 s
2 r23 1 2 r23 3
1 l 3 1 l 3
g1 s 1 s g3 s 3 s
6 r23 1 6 r23 3
This is because, by definition, e0 = l/r32. Thus, the quantity f1g3 − g1f3 can be

1 l 2 1 l 3 1 l 3 1 l 2
f1 g3 g1 f3 ¼ 1 s s3 s s1 s 1 s
2 r23 1 6 r23 3 6 r23 1 2 r23 3

1 l 3
¼ ðs3 s1 Þ s3 3s1 s23 þ 3s21 s3 s31
6 r23

1 l2 2 3
þ s1 s3 s31 s23
12 r26

1 l 3 1 l2 2 3
¼ ðs3 s1 Þ 3
ðs3 s1 Þ þ 6
s1 s3 s31 s23
6 r2 12 r2
Again, since s1 and s3 are small intervals, then the term (l2/r62) (s21s33 − s31s23)/12
can be neglected. Thus, setting s = s3 − s1 yields

1 l 3
f1 g3 g1 f3 s s
6 r23
where s is the interval of time elapsed from the first of the three observations to the
last. Now, substituting

1 l 3 1 l 3
f1 g3 g1 f3 s s g3 s 3 s
6 r23 6 r23 3
into c1 = g3/(f1 g3 − g1 f3) yields

1 l 3 1 l 2
s3 s 1 s

6 r23 3 s3 6 r23 3 s3 1 l 2 1 l 2 1
c1 ¼ ¼ ¼ 1 s 1 s
1 l 3 s 1 l 2 s 6 r23 3 6 r23
s s 1 s
6 r23 6 r23
Using the well-known binomial expansion
nðn 1Þ n2 2 nðn 1Þðn 2Þ n3 3

ða þ bÞn ¼ an þ nan1 b þ a b þ a b þ
2! 3!
1
with a ¼ 1; b ¼ ðl=r 32 Þs2 and n ¼ 1; truncated after the term containing the
6
second power of s, yields

1 l 2 1 1 1 l 2 1 l 2
1 s 1 þ ð1Þ1 11
s ¼ 1þ s
6 r23 6 r23 6 r23
This in turn substituted into

s3 1 l 2 1 l 2 1
c1 1 s 1 s
s 6 r23 3 6 r23
yields

s3 1 l 2 1 l 2 s3 1 l 2
c1 1 s 1þ s ¼ 1þ s s3
2
s 6 r23 3 6 r23 s 6 r23
By operating likewise, there results

s1 1 l 2
c3 1þ s s1
2
s 6 r23
We have hitherto obtained approximate expressions of the coefficients (c1 and

c3) appearing in the equation r2 = c1 r1 + c3 r3. These expressions depend only on
the (known) time intervals between observations and the (as yet unknown) distance
r2 of the attracted body from its centre of attraction at time t2.
The next stage of this development is the expression of the ranges q1, q2, and q3
as functions of c1 and c3. To this effect, we substitute
r1 ¼ rE1 þ q1 r1
r2 ¼ rE2 þ q2 u2
r3 ¼ rE3 þ q3 u3
into r2 = c1 r1 + c3 r3. This yields
rE2 þ q2 u2 ¼ c1 ðrE2 þ q1 u1 Þ þ c3 ðrE2 þ q3 u3 Þ
c1 q1 u1 q2 u2 þ c3 q3 u3 ¼ c1 rE1 þ rE2 c3 rE3
The scalar product of each term of the preceding expression by (u2 u3) yields
c1 q1 u1 ðu2 u3 Þ q2 u2 ðu2 u3 Þ þ c3 q3 u3 ðu2 u3 Þ

¼ c1 rE1 ðu2 u3 Þ þ rE2 ðu2 u3 Þ c3 rE3 ðu2 u3 Þ
Since u2 (u2 u3) = u3 (u2 u3) = 0, then the preceding expression

becomes
c1 q1 u1 ðu2 u3 Þ ¼ c1 rE1 ðu2 u3 Þ þ rE2 ðu2 u3 Þ c3 rE2 ðu2 u3 Þ
To simplify the notation, we set
D0 ¼ u1 ðu2 u3 Þ
D11 ¼ rE1 ðu2 u3 Þ
D21 ¼ rE2 ðu2 u3 Þ
D31 ¼ rE3 ðu2 u3 Þ
Assuming D0 6¼ 0, that is, assuming that u1, u2, and u3 are not coplanar, we
have
c1 q1 D0 ¼ c1 D11 þ D21 c3 D31
which, solved for q1, yields

1 c3 1
q1 ¼ D11 þ D21 D31
c1 c1 D0
By operating likewise, we take the scalar product of each term of
by (u1 u3). This yields


becomes
q2 u2 ðu1 u3 Þ ¼ c1 rE1 ðu1 u3 Þ þ rE2 ðu1 u3 Þ c3 rE3 ðu1 u3 Þ
Since −q2 u2 (u1 u3) = q2 u2 (u3 u1) = q2 u1 (u2 u3) = q2 D0, then
the preceding expression becomes
q2 D0 ¼ c1 rE1 ðu1 u3 Þ þ rE2 ðu1 u3 Þ c3 rE2 ðu1 u3 Þ

By setting
D12 ¼ rE1 ðu1 u3 Þ

D22 ¼ rE2 ðu1 u3 Þ
D32 ¼ rE2 ðu1 u3 Þ
and solving the preceding expression for q2, we obtain
1
q2 ¼ ðc1 D12 þ D22 c3 D32 Þ
D0
By operating likewise, we take the scalar product of each term of
by (u1 u2). This yields


becomes
c3 q3 u3 ðu1 u2 Þ ¼ c1 rE3 ðu1 u2 Þ þ rE3 ðu1 u2 Þ c3 rE3 ðu1 u2 Þ
By noting that u3 (u1 u2) = u1 (u2 u3) = D0, setting
D13 ¼ rE1 ðu1 u2 Þ

D23 ¼ rE2 ðu1 u2 Þ
D33 ¼ rE3 ðu1 u2 Þ
and solving the preceding expression for q3, we obtain

c1 1 1
q3 ¼ D13 þ D23 D33
c3 c3 D0
Now, substituting

s3 1 l s1 1 l
c1 1þ ðs s3 Þ
2 2
c3 1þ ðs s1 Þ
2 2
s 6 r23 s 6 r23
into q2 = (−c1 D12 + D22 − c3 D32)/D0 yields

l l
6þ ðs 2
s 2
3 Þ 6 þ ðs2 s21 Þ
r23 D22 r23
q2 ¼ D12 s3 þ þ D32 s1
6D0 s D0 6D0 s
which, after setting

s3 s1 1
A ¼ D12 þ D22 þ D32
s s D0

D12 s3 s2 s23 þ D32 s1 s2 s21
B¼
6D0 s
becomes
lB
q2 ¼ A þ
r23
Operating the same substitution, that is,

s3 1 l s1 1 l
c1 1þ ðs 2
s 2
3 Þ c 3 1 þ ðs 2
s 2
1 Þ
s 6 r23 s 6 r23
into
! !
1 c3 1 c1 1 1
q1 ¼ D11 þ D21 D31 0 q3 ¼ D13 þ D23 D33
c1 c1 D0 c3 c3 D0
leads to
2 s s1 3
6ðD31 ss13 þ D21 Þr23 þ lD31 ðs2 s21 Þ
6 s3 s3 7 1
q1 ¼ 4 D11 5
6r2 þ lðs s3 Þ
3 2 2 D0
2 s s3 3
6ðD13 ss31 D23 Þr23 þ lD13 ðs2 s23 Þ
6 s1 s1 7 1
q3 ¼ 4 D33 5
6r23 þ lðs2 s23 Þ D0
The equation written above q2 = A + lB/r32 expresses the range q2 as a function

of the radius vector r2 of the observed object measured from the centre of mass of
the Earth at time t2. Another relation between q2 and r2 is provided by
r2 ¼ rE2 þ q2 u2
The scalar product of r2 by itself yields
r2 r2 ¼ ðrE2 þ q2 u2 Þ ðrE2 þ q2 u2 Þ ¼ r 2E2 þ 2q2 ðrE2 u2 Þ þ q22 ¼ r 2E2 þ 2Eq2 þ q22
where E = rE2 u2. Substituting q2 = A + lB/r32 into r22 = r2E2 + 2Eq2 + q22 yields

lB lB 2
r22 ¼ 2
rE2 þ 2E A þ 3 þ A þ 3
r2 r2
2lEB 2lAB l2 B2
¼ rE2
2
þ 2EA þ þ A2
þ þ 6
r23 r23 r2
Multiplying all terms of the preceding expression by r62 yields

2
r28 rE2 þ 2EA þ A2 r26 2lBðE þ AÞr23 l2 B2 ¼ 0
If we set for convenience

x ¼ r2
2
a ¼ rE2 þ 2EA þ A2
b ¼ 2lB ðE þ AÞ
c ¼ l2 B2
the preceding expression becomes
x8 þ ax6 þ bx3 þ c ¼ 0
which is known as Lagrange’s equation. Since this polynomial has four terms, then
(according to the Descartes’ rule of signs, which states that the number of positive
roots of a polynomial with real coefficients is equal to the number of changes of
sign in the sequence of the coefficients of the polynomial, or is less than this number
by a multiple of 2) the polynomial may have no more than three positive roots.
When the correct root r2 of this equation has been found, it must be substituted into
2
s1 3
s1 s 3
6 D31 þ D21 r2 þ lD31 s2 s21
6 s3 s3 s3 7 1
q1 ¼ 6
4 D11
7
5 D0
6r2 þ lðs2 s3 Þ
3 2
lB
q2 ¼ A þ 3
r2
2 3
s3 s 3 s3
6 D13 D23 r2 þ lD13 ðs2 s23 Þ
6 s1 s1 s1 7 1
q3 ¼ 6
4 D33 7
5 D0
6r2 þ lðs s3 Þ
3 2 2
in order to compute the ranges q1, q2, and q3. Then, the following equations
r1 ¼ rE1 þ q1 u1
r2 ¼ rE2 þ q2 u2
r3 ¼ rE3 þ q3 u3
yield the position vectors r1, r2, and r3 of the observed object at times, respectively,
t1, t2, and t3. It remains to compute the velocity vector v2 of the observed object at
time t2. To this end, the equation r1 = f1 r2 + g1 v2 is solved for r2. This yields
r 1 g1 v 2
r2 ¼
f
The resulting value of r2 is substituted into r3 = f3 r2 + g3 v2. This yields
f3
r3 ¼ ðr1 g1 v2 Þ þ g3 v2
f1
and the preceding equation is solved for v2. This yields

f1 f3
v2 ¼ r3 r1
f1 g3 g1 f3 f1 g3 g1 f3
At first, the values of f1, f3, g1, and g3 to be used for computing r2 and v2 are
those derived previously, which are rewritten below for convenience:

1 l 2 1 l 2
f1 1 s f3 1 s
2 r23 1 2 r23 3

1 l 3 1 l 3
g1 s 1 s g3 s 3 s
6 r23 1 6 r23 3
Successively, improved values of fSince the condition1, f3, g1, and g3 are
computed, as will be shown below, and new values of r2 and v2 are computed by
means of these improved values, until convergence is reached. This method can be
illustrated by the following example, which is based on a series of astronomical
observations performed by Healy [33] and concerning the COBE artificial satellite
(USSPACECOM Catalogue No. 20322; International Designation code
1989-089-A). On the 6th of November 2000, Healy made seven observations of the
COBE satellite, three of which are shown in the following table.
Observed time (EST) Right ascension (hh:mm:ss) Declination (°)

17:31:29 21:48:00 −16.3
17:34:30 21:14:00 46.9
17:37:30 11:03:00 76.1
Since UTC = EST + 5, then (neglecting the difference between UTC and UT1)
the three EST times indicated above correspond, respectively, to
UT11 ¼ EST1 þ 5 ¼ 22h 31m 29s

UT12 ¼ EST2 þ 5 ¼ 22h 34m 30s
UT13 ¼ EST3 þ 5 ¼ 22h 37m 30s
The three values of the right ascension, expressed in degrees, are as follows
ð21 þ 48=60Þ 360=24 ¼ 327:00

ð21 þ 14=60Þ 360=24 ¼ 318:50
ð11 þ 3=60Þ 360=24 ¼ 165:75
Healy found these values by means of the 14″ Schmidt Cassegrain telescope of
the observatory of the University of Maryland, located at latitude u = 39°.00167
North, longitude k = −76°.95667 east and altitude H = 53 m.
As has been shown above, the same values can be written as follows
UT1 Right ascension (°) Declination (°)

22:31:29 327.00 −16.3
22:34:30 318.50 46.9
22:37:30 165.75 76.1
By applying the methods shown in Sect. 2.5, the Greenwich sidereal time hG0
corresponding to the 6 November 2000, at 00h:00m:00s UT1, results from
a ¼ INTðy=100Þ ¼ INTð2000=100Þ ¼ 20
b ¼ INTða=4Þ ¼ INTð20=4Þ ¼ 5
c ¼ 2 a þ b ¼ 2 20 þ 5 ¼ 13
e ¼ INT½365:25ðy þ 4716Þ ¼ INT½365:25 ð2000 þ 4716Þ ¼ 2453019
f ¼ INT½30:6001ðm þ 1Þ ¼ INT½30:6001 ð11 þ 1Þ ¼ 367
J0 ¼ c þ d þ e þ f 1524:5 ¼ 13 þ 6 þ 2453019 þ 367 1524:5 ¼ 2451854:5
T0 ¼ ðJ0 2451545Þ=36525 ¼ ð2451854:5 2451545Þ=36525 ¼ 0:0084736
hG0 ¼ 100:4606184 þ 36000:77005 T0 þ 0:000387933 T0 2 2:5875 108 T0 3
¼ 100:4606184 þ 36000:77005 0:0084736 þ 0:000387933 0:00847362
2:5875 108 0:00847363 ¼ 405 :51847
This value is brought into the range 0° hG0 360° by subtracting 360°, as
follows
hG0 ¼ 405 :51847 360 ¼ 45 :51847
The Greenwich sidereal times hG1, hG2, and hG3 corresponding to

UT11 = 22:31:29, UT12 = 22:34:30 and UT13 = 22:37:30 are, respectively,
hG1 ¼ hG0 þ 360:985647366 ð22 þ 31=60 þ 29=3600Þ=24 ¼ 384 :314363

hG2 ¼ hG0 þ 360:985647366 ð22 þ 34=60 þ 30=3600Þ=24 ¼ 385 :070595
hG3 ¼ hG0 þ 360:985647366 ð22 þ 37=60 þ 30=3600Þ=24 ¼ 385 :822648
The local sidereal times of the three observations are, respectively,
h1 ¼ hG1 þ k ¼ 384 :314363 76 :95667 ¼ 307 :357693

h2 ¼ hG1 þ k ¼ 385 :070595 76 :95667 ¼ 308 :113925
h3 ¼ hG1 þ k ¼ 385 :822648 76 :95667 ¼ 308 :865978
We want to compute the position and velocity vectors of the COBE satellite at
UT12 = 22h:34m:30s, with an accuracy of five significant figures. The angles
measured by the station and the local sidereal times are given in the following table.
Obs. No. Time (s) Right ascension (°) Declination (°) Local sidereal time (°)
1 0 327.00 −16.300 307.357693
2 181 318.50 46.900 308.113925
3 361 165.75 76.100 308.865978
First, we compute the three geocentric position vectors (rE1, rE2, and rE3) of the
observatory at the three given times t1, t2, and t3, as shown in Sect. 2.5 (that is,
taking the equatorial radius aE and the flattening f of the Earth equal, respectively,
to 6378.137 km and 0.0033528). Assuming the geodetic latitude equal to the
geographic latitude, the position vectors rE1, rE2, and rE3 are
n 1=2 o
xH ¼ aE = 1 2f f 2 sin2 u þ H cos u ¼ f6378:137=½1 ð2 0:0033528
o
0:00335282 Þ sin2 39:001671=2 þ 0:053 cos 39:00167 ¼ 4963:272 km
n 1=2 o n
zH ¼ aE ð1 f Þ2 = 1 2f f 2 sin2 u þ H sin u ¼ 6378:137 ð1 0:0033528Þ2
1=2 o
= 1 ð2 0:0033528 0:00335282 Þ sin2 39:00167 þ 0:053 sin 39:00167
¼ 3992:517 km
rE1 ¼ ðxH cos h1 Þ uX þ ðxH sin h1 Þ uY þ zH uZ ¼ ð4963:272 cos 307:357693Þ uX

þ ð4963:272 sin 307:357693Þ uY þ 3992:517 uZ
¼ 3011:7 uX 3945:1 uY þ 3992:5 uZ
rE2 ¼ ðxH cos h2 ÞuX þ ðxH sin h2 ÞuY þ zH uZ ¼ ð4963:272 cos 308:113925ÞuX
þ ð4963:272 sin 308:113925ÞuY þ 3992:517 uZ
¼ 3063:5 uX 3905:0 uY þ 3992:5 uZ
rE3 ¼ ðxH cos h3 ÞuX þ ðxH sin h3 ÞuY þ zH uZ ¼ ð4963:272 cos 308:865978ÞuX
þ ð4963:272 sin 308:865978ÞuY þ 3992:517 uZ
¼ 3114:5 uX 3864:5 uY þ 3992:5 uZ
The three unit vectors u1, u2, and u3 result from
u1 ¼ ðcos d1 cos a1 ÞuX þ ðcos d1 sin a1 ÞuY þ ðsin d1 ÞuZ ¼ ½cosð16:3Þ cos 327:0 uX
þ ½cosð16:3Þ sin 327:0uY þ ½sin ð16:3ÞuZ
¼ 0:80496 uX 0:52275 uY 0:28067 uZ
u2 ¼ ðcos d2 cos a2 Þ uX þ ðcos d2 sin a2 ÞuY þ ðsin d2 ÞuZ ¼ ðcos 46:9 cos 318:5ÞuX
þ ðcos 46:9 sin 318:5ÞuY þ ðsin 46:9ÞuZ
¼ 0:51174 uX 0:45275 uY þ 0:73016 uZ
u3 ¼ ðcos d3 cos a3 Þ uX þ ðcos d3 sin a3 ÞuY þ ðsin d3 ÞuZ ¼ ðcos 76:1 cos 165:75ÞuX
þ ðcos 76:1 sin 165:75ÞuY þ ðsin 76:1ÞuZ
¼ 0:23284 uX þ 0:059133 uY þ 0:97072 uZ
The time intervals s1, s3, and s are computed as follows
s1 ¼ t1 t2 ¼ 0 181 ¼ 181 s
s3 ¼ t3 t2 ¼ 361 181 ¼ 180 s
s ¼ t3 t1 ¼ 361 0 ¼ 361 s
The vector products (u2 u3), (u1 u3), and (u1 u2) result from
2 3
uX uY uZ
6 7
u2 u3 ¼ 4 0:51174 0:45275 0:73016 5
0:23284 0:059133 0:97072
¼ 0:48267 uX 0:66677 uY 0:075158 uZ
2 3
uX uY uZ
6 7
u1 u3 ¼ 4 0:80496 0:52274 0:28067 5
0:23284 0:059133 0:97072
¼ 0:49085 uX 0:71604 uY 0:074117 uZ
2 3
uX uY uZ
6 7
u1 u2 ¼ 4 0:80496 0:52274 0:28067 5
0:51174 0:45275 0:73016
¼ 0:50876 uX 0:73138 uY 0:096934 uZ
The scalar products D0, D11, D12, D13, D21, D22, D23, D31, D32, and D33 are
D0 ¼ u1 ðu2 u3 Þ ¼ 0:80496 ð0:48267Þ 0:522754 ð0:66677Þ

0:28067 ð0:075158Þ ¼ 0:018881
D11 ¼ rE1 ðu2 u3 Þ ¼ 3011:7 ð0:48267Þ 3945:1 ð0:66677Þ
þ 3992:5 ð0:075158Þ ¼ 876:75 km
D12 ¼ rE1 ðu1 u3 Þ ¼ 3011:7 ð0:49085Þ 3945:1 ð0:71604Þ
þ 3992:5 ð0:074117Þ ¼ 1050:6 km
D13 ¼ rE1 ðu1 u2 Þ ¼ 3011:7 ð0:50876Þ 3945:1 ð0:73138Þ
þ 3992:5 ð0:096934Þ ¼ 966:13 km
D21 ¼ rE2 ðu2 u3 Þ ¼ 3063:5 ð0:48267Þ 3905:0 ð0:66677Þ
þ 3992:5 ð0:075158Þ ¼ 825:01 km
D22 ¼ rE2 ðu1 u3 Þ ¼ 3063:5 ð0:49085Þ 3905:0 ð0:71604Þ
þ 3992:5 ð0:074117Þ ¼ 996:51 km
D23 ¼ rE2 ðu1 u2 Þ ¼ 3063:5 ð0:50876Þ 3905:0 ð0:73138Þ
þ 3992:5 ð0:096934Þ ¼ 910:44 km
D31 ¼ rE3 ðu2 u3 Þ ¼ 3114:5 ð0:48267Þ 3864:5 ð0:66677Þ
þ 3992:5 ð0:075158Þ ¼ 773:39 km
D32 ¼ rE3 ðu1 u3 Þ ¼ 3114:5 ð0:49085Þ 3864:5 ð0:71604Þ
þ 3992:5 ð0:074117Þ ¼ 942:47 km
D33 ¼ rE3 ðu1 u2 Þ ¼ 3114:5 ð0:50876Þ 3864:5 ð0:73138Þ
þ 4069:057 ð0:096934Þ ¼ 854:88 km
The quantities A and B are computed as follows
s3 =s ¼ 180=361 ¼ 0:49861
s1 =s ¼ 181=361 ¼ 0:50139

s3 s2 s23 ¼ 180 3612 1802 ¼ 1:7626 107
2
s1 s s21 ¼ 60 1202 602 ¼ 1:7658 107
A ¼ ðD12 s3 =s þ D22 þ D32 s1 =sÞ=D0 ¼ ½1050:6 0:49861 þ 996:51
þ 942:47 ð0:50139Þ=ð0:018881Þ ¼ 6:6363 km

B ¼ D12 s3 s2 s23 þ D32 s1 s2 s21 =½6D0 s ¼ 1050:6 1:7626 107
þ 942:47 ð1:7658 107 Þ=½6 ð0:018881Þ 361 ¼ 8:5974 108 km s2
The quantities E and r2E2 result from
E ¼ rE2 u2 ¼ 3063:5 0:51174 3905:0 ð0:45275Þ

þ 3992:5 0:73016 ¼ 6250:9 km
r 2E2 ¼ rE2 rE2 ¼ 3063:52 þ 3905:02 þ 3992:52 ¼ 4:0574 107 km2
The coefficients a, b, and c of the polynomial x8 + ax6 + bx3 + c result from

2 h i
a ¼ rE2 þ 2EA þ A2 ¼ 4:0574 107 þ 2 6250:9 ð6:6363Þ þ ð6:6363Þ2
¼ 4:0491 107 km2
b ¼ 2lBðE þ AÞ ¼ 2 398600:4 8:5974 108 ð6250:9 6:6363Þ
¼ 4:2797 1018 km5
2
c ¼ l2 B2 ¼ ð398;600:4Þ2 8:5974 108 ¼ 0:11744 1030 km8
We search a value of x such that the following function
f ðxÞ ¼ x8 þ ax6 þ bx3 þ c
should be equal to zero in a given interval xmin x xmax. This search is limited
to the values of x (where x r2) which are physically meaningful. Consequently,
x cannot be negative, or smaller than the mean radius of the Earth (6371 km). In
addition, since the values computed above of the three coefficients a, b, and c are,
all of them, negative, then the equation f(x) = x8 + ax6 + bx3 + c = 0 has only one
positive root. To search this root, we first evaluate f(x) for x ranging from 6400 km
to 9400 km, with a step size of 100 km. The results of this evaluation are shown
below.
x 103 (km) f(x) 1030 x 103 (km) f(x) 1030

6.400 −1.2071 6.900 −0.75509
6.500 −1.1601 7.000 −0.58431
6.600 −1.0941 7.100 −0.37856
6.700 −1.0067 7.200 −0.13376
6.800 −0.89472 7.300 0.15461
The plot shows that the physically meaningful root of f(x) = 0 is roughly the
midpoint of the interval 7200 x 7300 km. Therefore, the search of the root
of interest is confined to this interval. As shown in the preceding table, we find
f ð7200Þ ¼ 72008 4:0491 107 72006 4:2797 1018 72003

0:11744 1030 ¼ 0:13376 1030 ð\0Þ
f ð7300Þ ¼ 73008 4:0491 107 73006 4:2797 1018 73003
0:11744 1030 ¼ 0:15461 1030 ð [ 0Þ
Since the condition f(7200)f(7300) < 0 is satisfied, we search a zero of f(x) by

means of Müller’s method of parabolic interpolation, which has been described in
Chap. 1, Sects. 5 and 8. At the midpoint x0 = 7250 km of the interval we also find
f ð7250Þ ¼ 72508 4:0491 107 72406 4:2797 1018 72503

0:11744 1030 ¼ 0:00469 1030
Then we set
x2 ¼ 7200 km f2 f ðx2 Þ ¼ 0:13376 1030

x0 ¼ 7250 km f0 f ðx0 Þ ¼ 0:00469 1030
x1 ¼ 7300 km f1 f ðx1 Þ ¼ 0:15461 1030
and
h1 ¼ x1 x0 ¼ 7300 7250 ¼ 50 km
h2 ¼ x0 x2 ¼ 7250 7200 ¼ 50 km
c ¼ h2 =h1 ¼ 50=50 ¼ 1
and compute the coefficients

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½1 0:15461 1030 0:00469 1030 ð1
þ 1Þ 0:13376 1030 =½1 502 ð1 þ 1Þ ¼ 0:000002294 1030

B ¼ f1 f0 Ah21 =h1 ¼ ½0:15461 1030 0:00469 1030 0:000002294 1030
502 =50 ¼ 0:0028837 1030
C ¼ f0 ¼ 0:00469 1030
of the interpolating parabola f(x) = A(x − x0)2 + B(x − x0) + C.

We compute the estimated root of f(x) = 0, as follows
2C
x ¼ x0 1
B
ðB2 4 AC Þ2
where in the present case the plus sign in front of the square root takes effect,
because the value of B is greater than zero. Thus,
n
x ¼ 7250 2 0:00469 1030 = 0:0028837 1030 þ ½ð0:0028837 1030 Þ2
o
4 0:000002294 1030 0:00469 1030 1=2 ¼ 7248:4 km
This value, substituted into f(x) = x8 + ax6 + bx3 + c, yields
f ð7248:4Þ ¼ 7248:48 4:0491 107 7248:46 4:2797 1018 7248:43

0:11744 1030 ¼ 0:00009 1030
Now, since 7248.4 is less than 7250, then we take 7200, 7250 and 7248.4 for the
next step. At the same time, we reset the subscripts 0, 1 and 2, as follows
x2 ¼ 7200:0 km f2 f ðx2 Þ ¼ 0:13376 1030

x0 ¼ 7248:4 km f0 f ðx0 Þ ¼ 0:00009 1030
x1 ¼ 7250:0 km f1 f ðx1 Þ ¼ 0:00469 1030
Now we compute again
h1 ¼ x1 x0 ¼ 7250:0 7248:4 ¼ 1:6 km

h2 ¼ x0 x2 ¼ 7248:4 7200:0 ¼ 48:4 km
c ¼ h2 =h1 ¼ 48:4=1:6 ¼ 30:25
and the coefficients A, B, and C of the interpolating parabola, as follows

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½30:25 0:00469 1030 0:00009 1030
ð1 þ 30:25Þ 0:13376 1030 =½30:25 1:62 ð1 þ 30:25Þ ¼ 0:00000219 1030

B ¼ f1 f0 Ah21 =h1 ¼ ð0:00469 1030 0:00009 1030
0:00000219 1030 1:62 Þ=1:6 ¼ 0:0028715 1030
C ¼ f0 ¼ 0:00009 1030
Thus, the estimated root is

n
x ¼ 7248:4 2 0:00009 1030 = 0:0028715 1030 þ ½ð0:0028715 1030 Þ2
o
4 0:00000219 1030 0:00009 1030 1=2 ¼ 7248:4 km
Since this value is the same as that computed in the preceding step, then we take
7248.4 km as the correct root, within the chosen limits of accuracy, of
f ðxÞ ¼ x8 þ ax6 þ bx3 þ c ¼ 0
This means that we take x r2 = 7248.4 km in this preliminary approximation.

Now we compute the ranges q1, q2, and q3, as follows

2 3
s1 s 3 s1
6 D31 þ D21 r2 þ lD31 ðs2 s21 Þ
6 s3 s3 s3 7 1
q1 ¼ 6
4 D11 7
5 D0
6r2 þ lðs s3 Þ
3 2 2
lB
q2 ¼ A þ 3
r2
2 3
s3 s 3 s3
6 D13 D23 r2 þ lD13 ðs2 s23 Þ
6 s1 s1 s1 7 1
q3 ¼ 6
4 D33 7
5 D0
6r2 þ lðs s3 Þ
3 2 2

s1 =s3 ¼ 181=180 ¼ 1:0056
s=s3 ¼ 361=180 ¼ 2:0056
s=s1 ¼ 361=ð181Þ ¼ 1:9945
2
s s21 s1 =s3 ¼ ½3612 ð1812 Þ ð181=180Þ ¼ 98102
2
s s23 s3 =s1 ¼ ð3612 1802 Þ 180=ð181Þ ¼ 97380
s2 s23 ¼ 3612 1802 ¼ 97921
and consequently

q1 ¼ ½6 ð773:39 ð1:0056Þ þ 825:01 2:0056Þ 7248:43 þ 398600:4 773:39

ð98102Þ=½6 7248:43 þ 398600:4 97921 876:75 =ð0:018881Þ
¼ 1460:3 km
q2 ¼ 6:6363 þ 398600:4 8:5974 108 =7248:43 ¼ 893:23 km

q3 ¼ ½6 ð966:13=ð1:0056Þ 910:44 ð1:9945ÞÞ 7248:43 þ 398600:4 966:13

ð97380Þ=½6 7248:43 þ 398600:4 97921 854:88 =ð0:018881Þ
¼ 1602:4 km
We compute r1, r2, and r3 as follows
r1 ¼ rE1 þ q1 u1 ¼ 3011:7 uX 3945:1 uY þ 3992:5 uZ

þ 1460:3 ð0:80496 uX 0:52275 uY 0:28067 uZ Þ
¼ 4187:2 uX 4708:5 uY þ 3582:6 uZ
r2 ¼ rE2 þ q2 u2 ¼ 3063:5 uX 3905:0 uY þ 3992:5 uZ

þ 893:23 ð0:51174 uX 0:45275 uY þ 0:73016 uZ Þ
¼ 3520:6 uX 4309:4 uY þ 4644:7 uZ
r3 ¼ rE3 þ q3 u3 ¼ 3114:5 uX 3864:5 uY þ 3992:5 uZ

þ 1602:4 ð0:23284 uX þ 0:059133 uY þ 0:97072 uZ Þ
¼ 2741:4 uX 3769:7 uY þ 5548:0 uZ
The Lagrangian coefficients f1, f3, g1, and g3 result from
1 3 2 1
f1 1 l=r2 s1 ¼ 1 ð398600:4=7248:43 Þ ð181Þ2 ¼ 0:98285
2 2
1 3 2 1
f3 1 l=r2 s3 ¼ 1 ð398600:4=7248:43 Þ 1802 ¼ 0:98304
2 2
1 3 3 1
g1 s 1 l=r2 s1 ¼ 181 398600:4 ð181=7248:4Þ3 ¼ 179:97
6 6
1 3 3 1
g3 s 3 l=r2 s3 ¼ 180 398600:4 ð180=7248:4Þ3 ¼ 178:98
6 6
Now, in order to compute v2, we evaluate
f1 =ðf1 g3 g1 f3 Þ ¼ 0:98285=½0:98285 178:98 ð179:97Þ 0:98304 ¼ 2:7856 103

f3 =ðf1 g3 g1 f3 Þ ¼ 0:98304=½0:98285 178:98 ð179:97Þ 0:98304 ¼ 2:7862 103
v2 ¼ ½f1 =ðf1 g3 g1 f3 Þr3 ½f3 =ðf1 g3 g1 f3 Þr1 ¼ 0:0027856 ð2741:4 uX 3769:7 uY

þ 5548:0 uZ Þ 0:0027862 ð4187:2 uX 4708:5 uY þ 3582:6 uZ Þ
¼ 4:0299 uX þ 2:6179 uY þ 5:4727 uZ
In summary, the preliminary approximation to the position and velocity vectors

of the COBE satellite, observed at time t2, is
r2 ¼ 3520:6 uX 4309:4 uY þ 4644:7 uZ

v2 ¼ 4:0299 uX þ 2:6179 uY þ 5:4727 uZ
This completes the first part of the computation. The second part is meant to
compute more accurate values of the vectors r2 and v2 than those computed in the
first part, as will be shown below.
First iteration.
We compute the magnitudes r2 and v2 of the vectors r2 and v2, as follows
1
h i12
r2 ¼ ðr2 r2 Þ2 ¼ 3520:62 þ ð4309:4Þ2 þ 4644:72 ¼ 7248:4 km
1
h i12 1
v2 ¼ ðv2 v2 Þ2 ¼ ð4:0299Þ2 þ 2:61792 þ 5:47272 ¼ 53:0442 km=s
Now, we compute a, the inverse of the major semi-axis a of the trajectory of the
observed object, by means of the vis-viva integral v2/l = 2/r − 1/a, as follows
1 2 v22 2 53:044
a ¼ ¼ ¼ 1:4285 104 km1
a r2 l 7248:4 398600:4
We compute the radial component of v2, as follows
r2 r2 ð4:0299Þ 3520:6 þ 2:6179 ð4309:4Þ þ 5:4727 4644:7

v2r ¼ ¼
r2 7248:4
¼ 0:0069249 km/s
Then we write Kepler’s equation in universal variables (see Sect. 1.7) at times t1
and t3 respectively, as follows

1 r2 v2r 2 2
l ðt1 t2 Þ
2
1 v1 C av1 þ ð1 ar2 Þv31 S av21 þ r2 v1
l 2

1 r2 v2r 2 2
l ðt3 t2 Þ
2
1 v3 C av3 þ ð1 ar2 Þv33 S av23 þ r2 v3
l2
where C and S are the Stumpff functions (see Sect. 1.7), which are defined as
follows
8 1=2
< 1 cos z
> =z for z [ 0
C ðzÞ ¼ 1=2 for z ¼ 0
>
:
cosh z1=2 1 =ðzÞ for z\ 0
8 1=2 1=2 3=2
< z sin z
> =z for z [ 0
SðzÞ ¼ 1=6 for z ¼ 0
>
:
sinh z1=2 z1=2 =ðzÞ3=2 for z\ 0
z = av2, and v1 and v3 are the universal variables to be determined. Since
t1 t2 ¼ s1 ¼ 181
t3 t2 ¼ s3 ¼ 180
then Kepler’s equations at times t1 and t3 become, respectively,

h i
398600:41=2 ð181Þ ¼ 7248:4 ð0:0069249Þ=398600:41=2 v21 Cð1:4285 104 v21 Þ
þ ð1 1:4285 104 7248:4Þ v31 Sð1:4285 104 v21 Þ
þ 7248:4 v1
398600:41=2 180 ¼ ½7248:4 ð0:0069249Þ=398600:41=2 v23 Cð1:4285 104 v23 Þ

þ ð1 1:4285 104 7248:4Þv33 Sð1:4285 104 v23 Þ
þ 7248:4 v3
The preceding equations, after simplification, become, respectively,
114274 ¼ 0:079504 v21 Cð1:4285 104 v21 Þ 0:035434 v31 Sð1:4285 104 v21 Þ
þ 7248:4 vv1
113643 ¼ 0:079504 v23 Cð1:4285 104 v23 Þ 0:035434 v33 Sð1:4285 104 v21 Þ
þ 7248:4v3
We solve iteratively the two equations written above. To this end, an initial
estimate of v1 and v3 is provided by the following formula suggested by Chobotov
[18]:
1
l2 jajDt
As to the first equation, there results
l2 jajDt ¼ 398600:42 1:4285 104 ð181Þ ¼ 16:324 km2

1 1 1
Therefore, an estimate of the unknown value of v1 is taken tentatively in the

interval −18.0 v1 −14.0 around −16.324. We ascertain whether the fol-
lowing function
f ðvÞ ¼ 114274 þ 0:079504 v2 Cð1:4285 104 v2 Þ þ 0:035434 v3

Sð1:4285 104 v2 Þ 7248:4 v
has values of opposite signs at the endpoints of this interval. We find
f ð18:0Þ ¼ 114274 þ 0:079504 ð18:0Þ2 C½1:4285 104 ð18:0Þ2

þ 0:035434 ð18:0Þ3 S½1:4285 104 ð18:0Þ2 7248:4
ð18:0Þ ¼ 16176 ð [ 0Þ
f ð14:0Þ ¼ 114274 þ 0:079504 ð14:0Þ2 C½1:4285 104 ð14:0Þ2

þ 0:035434 ð14:0Þ3 S½1:4285 104 ð14:0Þ2 7248:4
ð14:0Þ ¼ 12805ð\0Þ
Since the condition f(−18.0)f(−14.0) < 0 is satisfied, we search a zero of f(v) by

means of Müller’s method of parabolic interpolation, which has been described in
Sects. 1.5 and 1.8. Consequently, we choose arbitrarily a value v0 falling between
the endpoints −18.0 and −14.0. By choosing v0 = −16.0, the corresponding f(v0) is
f ð16:0Þ ¼ 114274 þ 0:079504 ð16:0Þ2 C½1:4285 104 ð16:0Þ2

þ 0:035434 ð16:0Þ3 S½1:4285 104 ð16:0Þ2 7248:4 ð16:0Þ
¼ 1686:4
Then we set
v2 ¼ 18:0 km1=2 f2 f ðv2 Þ ¼ 16176

v0 ¼ 16:0 km1=2 f0 f ðv0 Þ ¼ 1686:4
v1 ¼ 14:0 km1=2 f1 f ðv1 Þ ¼ 12805
It is to be noted that here v2 and v1 are the endpoints of the current interval of
search for the unknown v and that they have nothing to do with the number of
observations. Now we set
h1 ¼ v1 v0 ¼ 14:0 ð16:0Þ ¼ 2:0 km1=2

h2 ¼ v0 v2 ¼ 16:0 ð18:0Þ ¼ 2:0 km1=2
c ¼ h2 =h1 ¼ 2:0=2:0 ¼ 1:0

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½1:0 ð12805Þ 1686:4 ð1
þ 1:0Þ þ 16176=½1:0 2:02 ð1 þ 1:0Þ ¼ 0:225

B ¼ f1 f0 A h21 =h1 ¼ ½12805 1686:4 ð0:225Þ 2:02 =2:0 ¼ 7245:3
C ¼ f0 ¼ 1686:4
of the interpolating parabola f(v) = A(v − v0)2 + B(v − v0) + C.

This done, we compute the estimated root of f(v) = 0 as follows
2C
v ¼ v0 1
B
ðB2 4AC Þ2
where in the present case the minus sign in front of the square root takes effect,
because the value of B is less than zero. Thus, there results
v ¼ 16:0 2 1686:4=½7245:3 ð7245:32 þ 4 0:225 1686:4Þ1=2

¼ 15:767 km1=2
This value, substituted into f(v), yields
f ð15:767Þ ¼ 114274 þ 0:079504 ð15:767Þ2 C ½1:4285 104 ð15:767Þ2

þ 0:035434 ð15:767Þ3 S ½1:4285 104 ð15:767Þ2 7248:4
ð15:767Þ ¼ 1:7312
Since −15.767 is greater than −16.0, we take −16.0, −14.0, and −15.767 for the
next step; at the same time, we reset the subscripts as follows
v2 ¼ 16:0 km1=2 f2 f ðv2 Þ ¼ 1686:4

v0 ¼ 15:767 km1=2 f0 f ðv0 Þ ¼ 1:7312
v1 ¼ 14:0 km1=2 f1 f ðv1 Þ ¼ 12805
Now we set
h1 ¼ v1 v0 ¼ 14:0 ð15:767Þ ¼ 1:767 km1=2

h2 ¼ v0 v2 ¼ 15:767 ð16:0Þ ¼ 0:233 km1=2
c ¼ h2 =h1 ¼ 0:233=1:767 ¼ 0:13186

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½0:13186 ð12805Þ ð1:7312Þ
ð1 þ 0:13186Þ þ 1686:4=½0:13186 1:7672 ð1 þ 0:13186Þ ¼ 0:23139

B ¼ f1 f0 Ah21 =h1 ¼ ½12805 ð1:7312Þ ð0:23139Þ 1:7672 =1:767
¼ 7245:4
C ¼ f0 ¼ 1:7312
of the interpolating parabola.

The estimated root v of f(v) = 0 results from
h 1=2 i
v ¼ v0 2C= B B2 4AC ¼ 15:767 2 ð1:7312Þ=½7245:4
ð7245:42 4 0:23139 1:7312Þ1=2 ¼ 15:767 km1=2
Since this value is the same as that computed previously, within the chosen
accuracy, we take it as correct. In other words, v1 = −15.767 km½ is taken as the
root of Kepler’s equation
114274 ¼ 0:079504v21 Cð1:4285 104 v21 Þ 0:035434 v31 Sð1:4285 104 v21 Þ
þ 7248:4v1
By using the same iterative method, we search now the root v3 of the second
equation, that is,
113643 ¼ 0:079504v23 Cð1:4285 104 v23 Þ 0:035434v33 Sð1:4285 104 v23 Þ

þ 7248:4v3
To this end, it is necessary to find an interval of v where the following function
f ðvÞ ¼ 113643 þ 0:079504v2 Cð1:4285 104 v2 Þ þ 0:035434v3 Sð1:4285 104 v2 Þ

7248:4v
changes sign. We try 14.0 v 18.0 and find the following values at the
endpoints:
f ð14:0Þ ¼ 113643 þ 0:079504 14:02 Cð1:4285 104 14:02 Þ þ 0:035434

14:03 Sð1:4285 104 14:02 Þ 7248:4 14:0 ¼ 12189 ð [ 0Þ
f ð18:0Þ ¼ 113643 þ 0:079504 18:02 Cð1:4285 104 18:02 Þ þ 0:035434

18:03 Sð1:4285 104 18:02 Þ 7248:4 18:0 ¼ 16781 ð\0Þ
Since the condition f(14.0)f(18.0) < 0 is satisfied, we choose arbitrarily a value

v0 falling between these endpoints. Choosing v0 = 16.0, the corresponding f(v0) is
f ð16:0Þ ¼ 113643 þ 0:079504 16:02 Cð1:4285 104 16:02 Þ þ 0:035434

16:03 Sð1:4285 104 16:02 Þ 7248:4 16:0 ¼ 2297:1
Now we set
v2 ¼ 14:0 km1=2 f2 f ðv2 Þ ¼ 12189

v0 ¼ 16:0 km1=2 f0 f ðv0 Þ ¼ 2297:1
v1 ¼ 18:0 km1=2 f1 f ðv1 Þ ¼ 16781
Again, the subscripts used here refer only to the iterations. We also set
h1 ¼ v1 v0 ¼ 18:0 16:0 ¼ 2:0 km1=2

h2 ¼ v0 v2 ¼ 16:0 14:0 ¼ 2:0 km1=2
c ¼ h2 =h1 ¼ 2:0=2:0 ¼ 1

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½1 ð16781Þ ð2297:1Þ
ð1 þ 1Þ þ 12189=½1 2:02 ð1 þ 1Þ ¼ 0:275

B ¼ f1 f0 Ah21 =h1 ¼ ½16781 ð2297:1Þ 0:275 2:02 =2:0 ¼ 7242:5
C ¼ f0 ¼ 2297:1
of the interpolating parabola f(v) = A(v − v0)2 + B(v − v0) + C.

Now, we compute the estimated root of f(v) = 0 as follows
h 1=2 i
v ¼ v0 2C= B B2 4AC ¼ 16:0 2 ð2297:1Þ=½7242:5 ð7242:52
þ 4 0:275 2297:1Þ1=2 ¼ 15:683
This value, substituted into f(v), yields
f ð15:683Þ ¼ 113643 þ 0:079504 15:6832 Cð1:4285 104 15:6832 Þ þ 0:035434

15:6833 Sð1:4285 104 15:6832 Þ 7248:4 15:683 ¼ 1:1684
Since 15.683 is less than 16.0, we take 14.0, 16.0 and 15.683 for the next
step. At the same time, the subscripts are reset as follows
v2 ¼ 14:0 km1=2 f2 f ðv2 Þ ¼ 12189

v0 ¼ 15:683 km1=2 f0 f ðv0 Þ ¼ 1:1684
v1 ¼ 16:0 km1=2 f1 f ðv1 Þ ¼ 2297:1
Now we set
h1 ¼ v1 v0 ¼ 16:0 15:683 ¼ 0:317 km1=2

h2 ¼ v0 v2 ¼ 15:683 14:0 ¼ 1:683 km1=2
c ¼ h2 =h1 ¼ 1:683=0:317 ¼ 5:3091

A ¼ ½cf1 f0 ð1 þ cÞ þ f2 = ch21 ð1 þ cÞ ¼ ½5:3091 ð2297:1Þ ð1:1684Þ
ð1 þ 5:3091Þ þ 12189=½5:3091 0:3172 ð1 þ 5:3091Þ ¼ 0:24895

B ¼ f1 f0 Ah21 =h1 ¼ ½2297:1 ð1:1684Þ 0:24895 0:3172 =0:317
¼ 7242:8
C ¼ f0 ¼ 1:1684
of the interpolating parabola.

The estimated root v of f(v) = 0 results from

h 1=2 i
v ¼ v0 2C= B B2 4AC ¼ 15:683 2 ð1:1684Þ=½7242:8
ð7242:82 þ 4 0:24895 1:1684Þ1=2 ¼ 15:683 km1=2
Since this value is the same as that computed previously, within the chosen
accuracy, we take it as correct. In other words, v3 = 15.683 km½ is taken as the root
of Kepler’s equation
113643 ¼ 0:079504v23 Cð1:4285 104 v23 Þ 0:035434v33 Sð1:4285 104 v23 Þ

þ 7248:4v3
Now, v1 = −15.767 km½ and v3 = 15.683 km½ are used to compute again the
Lagrangian coefficients f1, f3, g1 and g3 as follows
h i
f1 ¼ 1 ½v21 =r2 Cðav21 Þ ¼ 1 ð15:767Þ2 =7248:4 C ½1:4285 104 ð15:767Þ2
¼ 0:98290

f3 ¼ 1 ½v23 =r2 Cðav23 Þ ¼ 1 15:6832 =7248:4 Cð1:4285 104 15:6832 Þ
¼ 0:98308
g1 ¼ s1 ½v31 =l1=2 Sðav21 Þ ¼ 181 ½ð15:767Þ3 =398600:41=2
S½1:4285 104 ð15:767Þ2 ¼ 179:97 s
g3 ¼ s3 ½v33 =l1=2 Sðav23 Þ ¼ 180 ð15:6833 =398600:41=2 Þ
Sð1:4285 104 15:6832 Þ ¼ 178:98 s
These values of the Lagrangian coefficients are to be compared with those

resulting from the preliminary approximation. The two sets are shown below.
Preliminary-approximation values First-iteration values

f1 ¼ 0:98285 f1 ¼ 0:98290
f3 ¼ 0:98304 f3 ¼ 0:98308
g1 ¼ 179:97 g1 ¼ 179:97
g3 ¼ 178:98 g3 ¼ 178:98
In order to improve the convergence of the process, Curtis [20] suggests to

replace the first-iteration set of values by an arithmetic average of the two sets. This
leads to the following values
f1 ¼ ð0:98285 þ 0:98290Þ=2 ¼ 0:98288

f3 ¼ ð0:98304 þ 0:98308Þ=2 ¼ 0:98306
g1 ¼ 179:97
g3 ¼ 178:98
which are used to compute again c1 and c3 by means of

g3 g1
c1 ¼ c3 ¼
f1 g3 g1 f3 f1 g3 g1 f3
Thus,
c1 ¼ 178:98=½0:98288 178:98 ð179:97Þ 0:98306 ¼ 0:50726

c3 ¼ ð179:97Þ=½0:98288 178:98 ð179:97Þ 0:98306 ¼ 0:51007
These values of c1 and c3 are used to compute new values of q1, q2, and q3, as
follows
q1 ¼ ðD11 þ D21 =c1 c3 D31 =c1 Þ=D0 ¼ ð876:75 þ 825:01=0:50726

0:51007 773:39=0:50726Þ=ð0:018881Þ ¼ 1484:0 km
q2 ¼ ðc1 D12 þ D22 c3 D32 Þ=D0 ¼ ð0:50726 1050:6 þ 996:51
0:51007 942:47Þ=ð0:018881Þ ¼ 907:95 km
q3 ¼ ðc1 D13 =c3 þ D23 =c3 D33 Þ=D0 ¼ ð0:50726 966:13=0:51007
þ 910:44=0:51007 854:88Þ=ð0:018881Þ ¼ 1628:9 km
These values of q1, q2, and q3 are used to compute new values of r1, r2, and r3,
as follows
r1 ¼ rE1 þ q1 u1 ¼ 3011:7 uX 3945:1 uY þ 3992:5 uZ þ 1484:0 ð0:80496 uX

0:52275 uY 0:28067 uZ Þ ¼ 4206:3 uX 4720:9 uY þ 3576:0 uZ

0:45275 uY þ 0:73016 uZ Þ ¼ 3528:1 uX 4316:1 uY þ 4655:4 uZ

þ 0:059133 uY þ 0:97072 uZ Þ ¼ 2735:2 uX 3768:2 uY þ 5573:7 uZ
The values of r2 and r3 obtained above are used to compute a new value of v2, as
follows

f1 f3
v2 ¼ r3 r1
f1 g3 g1 f3 f1 g3 g1 f3
Since
f1 =ðf1 g3 g1 f3 Þ ¼ 0:98288=½0:98288 178:98 ð179:97Þ 0:98306 ¼ 0:0027856

f3 =ðf1 g3 g1 f3 Þ ¼ 0:98306=½0:98288 178:98 ð179:97Þ 0:98306 ¼ 0:0027862
then
v2 ¼ ½f1 =ðf1 g3 g1 f3 Þr3 ½f3 =ðf1 g3 g1 f3 Þr1 ¼ 0:0027856 ð2735:2 uX 3768:2 uY

þ 5573:7 uZ Þ 0:0027862 ð4206:3 uX 4720:9 uY þ 3576:0 uZ Þ
¼ 4:1004 uX þ 2:6567 uY þ 5:5626 uZ
In summary, the values of the position and velocity vectors of the COBE satellite
at time t2, at this stage of the iterative procedure, are
r2 ¼ 3528:1 uX 4316:1 uY þ 4655:4 uZ

v2 ¼ 4:1004 uX þ 2:6567 uY þ 5:5626 uZ
The remaining part of the computation is straightforward, but long. Therefore, it

is not given here. Suffice it to say that, after the fifth iteration, the values of the four
Lagrangian coefficients converge to those given below:
f1 ¼ 0:98297
f3 ¼ 0:98315
g1 ¼ 179:97
g3 ¼ 178:99
Consequently, the iterative process of refinement terminates after the fifth iter-
ation. The position (r2) and velocity (v2) vectors of the COBE satellite at time t2 can
be computed by introducing these values in the expressions shown above. The
computed values of these vectors are given below.
r2 ¼ 3525:5 uX 4313:7 uY þ 4651:7 uZ

v2 ¼ 4:0755 uX þ 2:6425 uY þ 5:5324 uZ
The computed values can be compared with the actual values of the position and
velocity vectors of the COBE satellite, at the same time t2, which are given below
(from Healy [33]).
r2 ¼ 3528:320 uX 4313:871 uY þ 4654:938 uZ

v2 ¼ 4:103 uX þ 2:658 uY þ 5:564 uZ
About the difference existing between computed and actual values, it is to be

noted that the data used for the computation have been considered as free from
errors, which cannot be true in practice; in addition, no account has been taken of
the perturbations.
The method shown above is due to Gauss [26] and has been refined by Gibbs
[27]. Taff [66] has strongly censured this method on the grounds of the small radius
of convergence of the f and g series. According to Taff, Moulton [56] has found that
the radius of convergence of these series is Ts/2p, where T is the orbital period (of
the object observed) and s is given by
n h 1 i 1 o2
s2 ¼ M02 þ ln 1 þ 1 e2 2 ln e 1 e2 2
where e is the orbital eccentricity and M0 is the value of the mean anomaly,
contained in the interval [−p, p], at the instant t = t0 [66]. Neither Taff nor Moulton,
in the cited papers, define what must be understood by “instant t = t0”. It may be
presumed either that t0 is the intermediate (t2) of the three observation times (t1, t2,
and t3) defined above, which are such that t1 < t2 < t3, or that t0 is exactly the
midpoint of the interval [t1, t3], that is, t0 = (t1 + t3)/2.
On the other hand, well before Taff and Moulton, Gibbs had noted that “The
determination of an orbit from three complete observations by the solution of the
equations which represent elliptic motion present so great difficulties in the general
case, that in the first solution of the problem we must generally limit ourselves to
the case in which the intervals between the observations are not very long” [27]. In
other words, the time intervals s1 = t1 − t2 and s3 = t3 − t2 must be small fractions
of the period T of the orbiting body observed, which period is not known a priori.
However, the method of Gauss shown above computes first the components of the
position and velocity vectors of the observed object from a preliminary approxi-
mation and then refines iteratively the values of such components. The knowledge
of the first-approximation set of values makes it possible to compute approximately
the major semi-axis a = 1/(2/r − v2/l) of the orbit and hence the orbital period
T = 2p(a3/l)½ of the observed object.
In addition, Branham [13] notes that, according some authors, Gauss’ method is
restricted to low-eccentricity orbits, because the radius of convergence of the f and
g series becomes smaller and smaller as the orbital eccentricity approaches unity.
This is because the value of s in Moulton’s [56] formula given above also depends
on eccentricity. However, Branham himself admits that this restriction can be
removed by using the f and g functions rather than series.
On the same line of reasoning is Marsden, who shows that Taff’s criticism of
Gauss’ method of initial orbit determination on the grounds of the small radius of
convergence of the f and g series is completely unjustified [51].
In the example given above, the preliminary approximation to the position and
velocity vectors of the COBE satellite, observed at time t2, has given the following
results
r2 ¼ 3520:6 uX 4309:4 uY þ 4644:7 uZ

v2 ¼ 4:0299 uX þ 2:6179 uY þ 5:4727 uZ
The corresponding magnitudes r2 and v2 of r2 and v2 have been found to be
1
h i12
r2 ¼ ðr2 r2 Þ2 ¼ 3520:62 þ ð4309:4Þ2 þ 4644:72 ¼ 7248:4 km
1
h i12 1
v2 ¼ ðv2 v2 Þ2 ¼ ð4:0299Þ2 þ 2:61792 þ 5:47272 ¼ 53:0442 km=s
The major semi-axis results from the vis-viva integral
1 1
a¼ ¼ ¼ 7000:5 km
2 v22 2 53:044

r2 l 7248:2 398600:4
The corresponding orbital period of the COBE satellite results from

3 12 1
a 7000:53 2
T ¼ 2p ¼ 2 3:1416 ¼ 5829:1 s
l 398600:4
and the angular momentum per unit mass is

h ¼ Y2 Z20 Z2 Y20 uX þ Z2 X20 X2 Z20 uY þ X2 Y20 Y2 X20 uZ
¼ ð4309:4 5:4727 2:6179 4644:7Þ uX þ ð3520:6 5:4727
4:0299 4644:7Þ uY þ ð3520:6 2:6179 4:0299 4309:4Þ uZ
¼ 35743 uX 37985 uY 8149:9 uZ
The square of the magnitude of h is
h2 ¼ h h ¼ ð35743Þ2 þ ð37985Þ2 þ ð8149:9Þ2
The semi-latus rectum results from
h2 ð35743Þ2 þ ð37985Þ2 þ ð8149:9Þ2

p¼ ¼ ¼ 6991:6 km
l 398600:4
The orbital eccentricity results from

p ¼ a 1 e2
which, solved for e, yields
1
p12 6991:6 2
e¼ 1 ¼ 1 ¼ 0:035656
a 7000:5
Now, since in the present case

• the time intervals s1 = −181 s and s3 = 180 s are equal to a small fraction
(about 1/32) of the orbital period T = 5829.1 s; and
• the orbital eccentricity e = 0.035656 is much less than unity;
then the iterations converge, as has been shown above.

(Method of Laplace)
Let the topocentric right ascension (a) and declination (d) of a satellite be given at
three distinct times t1, t2, and t3, that is, let a set of values for a1, d1, a2, d2, a3, and
d3 be given. Let u1, u2, and u3 be the unit vectors along the line-of-sight vectors
(respectively, q1 q1u1, q2 q2u2, and q3 q3u3) going from the observation
site P placed on the surface of the Earth to the satellite Q observed, as shown in the
following figure.
Indicating for brevity ui (i = 1, 2, and 3) these unit vectors, Sect. 2.2 has shown
that the components of ui are
ui ¼ ðcos di cos ai ÞuX þ ðcos di sin ai ÞuY þ ðsin di ÞuZ

2.7 Orbital Elements from Three Measurements of Angles (Method of Laplace) 213
where uX, uY, and uZ are the unit vectors along the respective axes of the
geocentric-equatorial system XYZ. The preceding expression, in matrix terms, is
2 3 2 3
uX cos d cos a
ui ¼ 4 uY 5 ¼ 4 cos d sin a 5 ði ¼ 1; 2; 3Þ
uZ i sin d i
The preceding figure also shows that
ri ¼ rE þ qi ¼ rE þ qi ui ði ¼ 1; 2; 3Þ
where qi is the magnitude of the line-of-sight vector qi, and rE is the vector from the
centre of the Earth to the observation site.
As shown in Sect. 2.5, differentiating two times the preceding expression with
respect to time yields
r0i ¼ r0E þ q0i ui þ qi u0i

r00i ¼ r00E þ q00i ui þ q0i u0i þ q0i u0i þ qi u00i ¼ r00E þ q00i ui þ 2q0i u0i þ qi u00i
On the other hand, the equation of motion of the satellite is

l
r00i ¼ 3 ri
ri
By substituting r″i = −(l/r3i )ri into r″i = r″E + q″i ui + 2 q′i u′i + qi u″i and
remembering that ri = rE + qi ui, there results

l
ðrE þ qi ui Þ ¼ r00E þ q00i ui þ 2q0i u0i þ qi u00i
ri3
that is,
h l i h l i
q00 ui þ 2q0i u0i þ qi u00i þ 3 ui ¼ r00E þ 3 rE
r r
(where i = 1, 2, 3). At a given time of observation (e.g. at the intermediate time t2)
the preceding vector equation is equivalent to 3 scalar equations in 10 unknowns
(u′i, u″i, qi, q′i, q″i, and r), where the known quantities are ui, rE, and r″E. The
line-of-sight vector ui is known at the three times t1, t2, and t3. Consequently, when
the time intervals t2 − t1 and t3 − t2 are small with respect to the orbital period T of
the object observed, the values of u′2 and u′′2 (that is, the values of u′i and u′′i at the
intermediate time t2) can be computed by taking, respectively, the first and second
time-derivative of u(t), where u(t) is the Lagrange polynomial which interpolates
the three line-of-sight vectors u1, u2, and u3. As is well known, this polynomial is

ðt t 2 Þðt t 3 Þ ðt t 1 Þðt t 3 Þ ðt t1 Þðt t2 Þ
uðtÞ ¼ u1 þ u2 þ u3
ðt 1 t 2 Þðt 1 t 3 Þ ðt2 t1 Þðt2 t3 Þ ðt3 t1 Þðt3 t2 Þ
Differentiating once and twice this expression yields, respectively,

2t t2 t3 2t t1 t3 2t t1 t2
u0 ðtÞ ¼ u1 þ u2 þ u3
ðt1 t2 Þðt1 t3 Þ ðt2 t1 Þðt2 t3 Þ ðt3 t1 Þðt3 t2 Þ

2 2 2
u00 ðtÞ ¼ u1 þ u2 þ u3
ðt1 t2 Þðt1 t3 Þ ðt2 t1 Þðt2 t3 Þ ðt3 t1 Þðt3 t2 Þ
The value of u′2 results from evaluating u′(t2). When more than three obser-
vations are available, then u′2 and u′′2 can be computed by using either Lagrange
polynomials with a higher degree than two or a least-squares polynomial fit, as will
be shown below.
Thus, at t = t2, we have the equation
h l i h l i
q00 u þ 2q0 u0 þ q u00 þ 3 u ¼ r00E þ 3 rE
r r
which, projected onto the axes X, Y, and Z, gives rise to the following scalar
equations
q00 uX þ 2q0 u0X þ q½u00X þ ðl=r 3 Þ‘X ¼ ½rEX

00
þ ðl=r 3 ÞrEX
q00 uY þ 2q0 u0Y þ q½u00Y þ ðl=r 3 Þ‘Y ¼ ½rEY
00
þ ðl=r 3 ÞrEY
q00 uZ þ 2q0 u0Z þ q½u00Z þ ðl=r 3 Þ‘Z ¼ ½rEZ
00
þ ðl=r 3 ÞrEZ
for the three unknowns q″, q′, and q. For the present, we consider r as a parameter.
The matrix of the coefficients in the preceding system of equations is
2 3
uX 2u0X u00X þ ðl=r 3 ÞuX
4 uY 2u0Y u00Y þ ðl=r 3 ÞuY 5
uZ 2u0Z u00Z þ ðl=r 3 ÞuZ
Let D be the determinant of the preceding matrix. The value of D does not
change if the first column of the matrix, multiplied by (l/r3), is subtracted from the
third column, as follows
2 3
uX 2u0X u00X þ ðl=r 3 ÞuX uX ðl=r 3 Þ
4 uY 2u0Y u00Y þ ðl=r 3 ÞuY uY ðl=r 3 Þ 5
uZ 2u0Z u00Z þ ðl=r 3 ÞuZ uZ ðl=r 3 Þ
Consequently, D is equal to
2 3
uX u0X u00X
24 uY u0Y u00Y 5 ¼ 2D0
uZ u0Z u00Z
where D0 is the determinant of the matrix written above. Solving the preceding
system of equations according to Cramer’s rule, we have
Dq Dq
q¼ ¼
D 2D0
where Dq is the determinant of the following matrix

2 3
uX 2u0X 00
rEX þ ðl=r 3 ÞrEX
4 uY 2u0Y 00
rEY þ ðl=r 3 ÞrEY 5
uZ 2u0Z 00
rEZ þ ðl=r 3 ÞrEZ
This matrix can be written as follows

2 3 2 3
uX u0X 00
rEX uX u0X rEX
24 uY u0Y 00 5
rEY 2 l=r 3 4 uY u0Y rEY 5
uZ u0Z 00
rEZ uZ u0Z rEZ
Let D1 and D2 be the determinants of the two matrices written above.

Since q = Dq/D, then the following equality holds
Dq D1 lD D1 l D2 l
2
q¼ ¼ 2 2 3 ¼ 3 ¼ A B 3
D D r D D0 r D0 r
where we have set for convenience

D1 u u0 r00E
A¼ ¼
D0 D0
D2 u ðu0 rE Þ
B¼ ¼
D0 D0
0 00
D 0 ¼ u ðu u Þ
The equation written above, that is,

l
q ¼ A B 3
r
holds, of course, if D0 6¼ 0.
On the other hand, by remembering that r = rE + q u (where r and u are ri and ui

at the intermediate time t2) and taking the scalar product of r with itself, there results
r2 ¼ r r ¼ ðrE þ q uÞ ðrE þ q uÞ ¼ r 2E þ 2 q ðrE uÞ þ q2
By setting for convenience
E ¼ 2 q ðrE uÞ
F ¼ rE2
there results the following system of two algebraic equations

l
q ¼ A B 3
r
r 2 ¼ q2 E q þ F
for the two unknowns q and r.

Now, substituting q = −A − B(l/r3) into r2 = q2 − E q + F leads to
2l A B l2 B2 l BE
r 2 ¼ A2 þ þ 6 þ AE þ 6 þ F
r3 r r
Multiplying all terms of the preceding equation by r6 leads to Lagrange’s equation
r 8 þ ar 6 þ br 3 þ c ¼ 0
where

a ¼ A2 þ AE þ F
b ¼ lð2AB þ BE Þ
c ¼ l2 B2
This equation can be solved numerically by means of one of the methods shown
in Sect. 1.5. The result is the value of r relating to the intermediate time t2. By
introducing this value into
l
q ¼ A B 3
r
it is possible to compute q. Thus, the position vector of the satellite at time t2 results
from
r ¼ rE þ q u
With the view of computing the velocity vector v r′ of the satellite at time t2,
let us consider again the system of equations
00
q00 uX þ 2q0 u0X þ q u00X þ ðl=r 3 Þ‘X ¼ rEX þ ðl=r 3 ÞrEX
00
q00 uY þ 2q0 u0Y þ q u00Y þ ðl=r 3 Þ‘Y ¼ rEY þ ðl=r 3 ÞrEY
00
q00 uZ þ 2q0 u0Z þ q u00Z þ l=r 3 ‘Z ¼ rEZ þ l=r 3 rEZ
Using again Cramer’s rule, we solve now for q′ and write
D q0
q0 ¼
D
where D 6¼ 0 is the determinant of the matrix of coefficients (see above) and Dq′ is
the determinant of the following matrix
2 00 3
uX rEX þ ðl=r 3 ÞrEX u00X þ ðl=r 3 ÞuX
4 uY 00
rEY þ ðl=r 3 ÞrEY u00Y þ ðl=r 3 ÞuY 5
00
uZ rEZ þ ðl=r 3 ÞrEZ u00Z þ ðl=r 3 ÞuZ
This matrix can also be written as follows

2 00
3 2 3
uX rEX u00X uX rEX u00X
4 uY 00
rEY u00Y 5 l=r 3 4 uY rEY u00Y 5
00
uZ rEZ u00Z uZ rEZ u00Z
Let D3 be the determinant of the first and D4 be the determinant of the second of
the two matrices written above. By so doing, the previous expression q′ = Dq′/D
D q0 D3 l D4
q0 ¼ ¼ 3
D D r D
where D 6¼ 0. Since we have set D = 2D0, the previous equation can also be written
as follows
D3 l D
q0 ¼
4
3
2D0 r 2D0
By setting for convenience

D3 u r00E u00
C¼ ¼
2D0 2D0
D4 u ðrE u00 Þ
G¼ ¼
2D0 2D0
the previous equation can be written as follows

l
q0 ¼ C G 3
r
Since r is known, the previous expression makes it possible to compute the value
of q′. This value, substituted into
v r0 ¼ r0E þ q0 u þ q u0
yields the velocity vector v of the satellite at time t2.

As mentioned above, this method fails when the determinant D0 of the matrix
2 3
uX u0 X u00 X
4 uY u0 Y u00 Y 5
uZ u0 Z u00 Z
approaches zero, which happens when the observer lies in the plane of the orbit at
the intermediate time t2.
The following section of this paragraph shows how to fit the observations
gathered to an orbit by using the method of least squares, which was invented by
Gauss in 1795. According to Gauss (Theoria motus, Book 2, Sect. 3, paragraph
172, page 205 of the cited Ref. [26]), “Si observationes astronomicae ceterique
numeri, quibus orbitarum computus innititur, absoluta praecisione gauderent, ele-
menta quoque, sive tribus observationibus sive quatuor superstructa fuerint, abso-
lute exacta statim prodirent (quatenus quidem motus secundum leges Kepleri exacte
fieri supponitur), adeoque accitis aliis aliisque observationibus confirmari tantum
possent, haud corrigi. Verum enim vero quum omnes mensurationes atque obser-
vationes nostrae nihil sint nisi approximationes ad veritatem, idemque de omnibus
calculis illis innitentibus valere debent, scopum summum omnium computorum
circa phaenomena concreta institutorum in eo ponere oportebit, ut ad veritatem
quam proxime fieri potest accedamus. Hoc autem aliter fieri nequit, nisi per ido-
neam combinationem observationum plurium, quam quot ad determinationem
quantitatum incognitarum absolute requiruntur. Hoc negotium tunc demum susci-
pere licebit, quando orbitae cognitio approximata iam innotuit, quae dein ita rec-
tificanda est, ut omnibus observationibus quam exactissime satisfaciat”.
In plain English, “If the astronomical observations and other numbers, on which
the computation of orbit is based, were absolutely precise, the elements also,
deduced by means of three or four observations, would be absolutely exact (within
the limits of validity of Kepler’s laws). Hence, the computed elements could only
be confirmed, but never corrected, by further observations. In practice, since all our
measurements and observations are nothing more than approximations to the truth,
the same holds for all calculations based on them. Thus, the principal purpose of all
computations concerning concrete phenomena must needs be that of drawing us
near the truth to the maximum possible extent. This can be done in no other way
than by using a suitable combination of more observations than would be strictly

necessary to determine the unknown quantities. This task can only be undertaken,
when an approximate knowledge of the orbit has already been reached, which is
then to be corrected in such a way as to satisfy all the observations to the maximum
extent of exactness”.
In a few words, since all actual observations are affected by random errors, then
more observations than those which are strictly necessary are needed, the usefulness
of the redundant measurements being the mutual cancellation of the errors to the
maximum possible extent.
With reference to the following figure, suppose we have a set of approximate
values yi which correspond to discrete values xi (where i = 1, 2, …, n) of the
independent variable x.
If the three points (x1, y1), (x2, y2), and (x3, y3) were joined by an interpolating
polynomial, as shown in the following figure, then the resulting curve would pass
exactly through each of the three points but would also have an oscillating beha-
viour, which makes the curve unfit to represent the overall trend of our data.
If, on the contrary, the same data were represented by a straight line obtained by
means of the least-squares method, then the result would be like that shown below.
The least-squares method produces a functional form f(x, c0, c1, …, cm) fm(x),
which depends not only on x but also on m + 1 parameters c0, c1, …, cm to be
determined in such a way as to minimise the squares of the residuals, that is, the
squares of the differences between the functional form fm(x) and each data point.
A type of functional form fm(x) frequently used is an algebraic polynomial:
f m ð xÞ ¼ c0 þ c1 x þ c2 x2 þ þ cm xm
The degree m of this polynomial must be chosen by the solver. However, m must
be less than n − 1, where n is the number of the given points. Otherwise, if m were
equal to n − 1, then fm(x) would be just the interpolating polynomial.
For the purpose of illustrating the method, suppose that we are given a set of
n points (xi, yi) with i = 1, 2, …, n, and that we search the least-squares fit of these
points by means of a second-degree polynomial f2(x) = c0 + c1x + c2x2. In this
case, according to the least-squares method, the unknown coefficients c0, c1, and c2
must be chosen in such a way as to correspond to the least squares of the residuals,
the residuals being the components q1, q2, …, qn of the following n 1 residual
vector
23
q1
6 q2 7
6 7
q6 . 7
4 .. 5
qn
Such components are
q1 ¼ y1 f 2 ðx1 Þ ¼ y1 c0 c1 x1 c2 x21
q2 ¼ y2 f 2 ðx2 Þ ¼ y2 c0 c1 x2 c2 x22
..
.
qn ¼ yn f 2 ðxn Þ ¼ yn c0 c1 xn c2 x2n
This is because the Euclidean length q of the residual vector q is just the square
root of the sum of the squares of its components: q = (q21 + q22 + … + q2n)½.
Consequently, in order for the residual vector q to have the minimum Euclidean
length, q2 must have the minimum value, where
2
q2 ¼ q21 þ q22 þ þ q2n ¼ y1 c0 c1 x1 c2 x21
2 2
þ y2 c0 c1 x2 c2 x22 þ þ yn c0 c1 xn c2 x2n
This condition is satisfied by those values of c0, c1, …, cm which cause the first
derivative of q2 to vanish. If m = 2, this leads to the following three conditions:
@q2 @q21 @q22 @q2

¼ þ þ þ n ¼ 2½ y1 c0 c1 x1 c2 x21
@c0 @c0 @c0 @c0

þ y2 c0 c1 x2 c2 x22 þ þ yn c0 c1 xn c2 x2n ¼ 0
@q2 @q21 @q22 @q2
¼ þ þ þ n ¼ 2½x1 y1 c0 c1 x1 c2 x21
@c1 @c1 @c1 @c1

þ x2 y2 c0 c1 x2 c2 x22 þ þ xn yn c0 c1 xn c2 x2n ¼ 0
@q2 @q21 @q22 @q2
¼ þ þ þ n ¼ 2½x21 y1 c0 c1 x1 c2 x21
@c2 @c2 @c2 @c2

þ x22 y2 c0 c1 x2 c2 x22 þ þ x2n yn c0 c1 xn c2 x2n ¼ 0
Since each of the three partial derivatives must be equal to zero, then the
coefficient in front of the square brackets (that is, −2) can be dropped. Thus, the
previous three conditions reduce to

nc0 þ ðx1 þ x2 þ . . . þ xn Þc1 þ x21 þ x22 þ . . . þ x2n c2 ¼ y1 þ y2 þ . . . þ yn

ðx1 þ x2 þ . . . þ xn Þc0 þ x21 þ x22 þ . . . þ x2n c1 þ x31 þ x32 þ . . . þ x3n c2
¼ x1 ðy1 þ y2 þ . . . þ yn Þ
2
x1 þ x22 þ . . . þ x2n c0 þ x31 þ x32 þ . . . þ x3n c1 þ x41 þ x42 þ . . . þ x4n c2
¼ x21 ðy1 þ y2 þ . . . þ yn Þ
This is a system of three linear algebraic equations for the three unknowns c0, c1,
and c2. The values of c0, c1, and c2 determine the least-squares fit parabola
f 2 ð xÞ ¼ c0 þ c1 x þ c2 x2
Let us consider now the differentiation of a general (m-degree) least-squares fit

polynomial. Let
fm ðx ¼ c0 þ c1 x þ c2 x2 þ . . . þ cn xm
be this polynomial. Differentiating once and twice the polynomial with respect to
x yields
dfm
¼ c1 þ 2c2 x þ 3c3 x2 þ
dx
d 2 fm
¼ 2c2 þ 6c3 x þ
dx2
The expressions shown above make it possible to compute fm(x) and its first and
second derivative at any point x of interest. If that point be chosen as the origin
of the series expansions, then there results x = 0, fm(0) = c0, (dfm/dx)0 = c1, and
(d2fm/dx2)0 = 2c2.
2.8 Improvement in Orbit Determination by Differential

Correction
The preceding paragraphs have shown some methods for the preliminary deter-
mination of the orbit followed by a space object. These methods use only the
minimum number of observations necessary to compute an orbit. Since six inde-
pendent quantities are strictly necessary to determine the motion of a body, then
these methods use a set of six independent quantities.
However, as has been shown numerically in an example given in Sect. 2.6, the
six orbital elements (or the 3 + 3 components of the position and velocity vectors at
a given time t0) computed by means of these methods differ to some extent from the
actual orbital elements. For the purpose of reducing, as far as possible, the influence
of the errors, the residuals (that is, the differences between the computed and
observed quantities) must be determined at each time of observation. In some cases,
it is possible to refine a preliminary orbit without taking account of the perturba-
tions, that is, considering a purely Keplerian orbit. As shown in Sect. 2.6, this
happens when there are at least three reliable observations which cover a significant
part of the orbit but are not too distant in time with respect to the epoch of the
preliminary determination of the orbital elements.
The differential correction is a numerical procedure based on the principle of least
squares, which is meant to correct the computed elements by minimising the
residuals. The quantities which concur to determine the observed values of the right
ascension (a) and declination (d) of a given body are substantially the orbital ele-
ments of that body and the position vector rE which specifies the location of the
observer with respect to the geocentric-equatorial system. Of these quantities, rE can
be supposed to be known accurately enough to require no improvement in accuracy
make it useless to improve rE. Consequently, the uncertainties in the values of a and
d can be attributed to the orbital elements, that is, to the components of the position
(r) and velocity (r′) vectors of the space object. Let t0 be an epoch chosen arbitrarily
and let x0, y0, z0, x′0, y′0, and z′0 be the components of such vectors at t0. The
dependency of a and d on x0, y0, z0, x′0, y′0, and z′0 can be written as follows

a ¼ a x0 ; y0 ; z0 ; x00 ; y00 ; z00

d ¼ d x0 ; y0 ; z0 ; x00 ; y00 ; z00
The total differentials of a and d are
@a @a @a
da ¼ dx0 þ dy0 þ þ 0 dz00
@x0 @y0 @z0
@d @d @d
dd ¼ dx0 þ dy0 þ þ 0 dz00
@x0 @y0 @z0
and express the changes in a and d resulting from the independent changes in x0, y0,
…, z′0. Replacing the differentials by the corresponding finite differences leads to
2.8 Improvement in Orbit Determination by Differential Correction 223
@a @a @a
Da ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
@d @d @d
Dd ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
where Da and Dd are the residuals in, respectively, right ascension and declination,
and Dx0, Dy0, …, Dz′0 are the changes to be made in, respectively, x0, y0, …, z′0 for
the purpose of reducing such residuals to zero. Now, Da and Dd are known
quantities, which result from measuring the position of the observed body on the
celestial sphere. When three couples of values (a, d) are measured at three different
times t1, t2, and t3, that is, when we have a set of six values {a1, d1, a2, d2, a3, d3},
then 2 3 independent equations of conditions can be written as follows
@a1 @a1 @a1 0

Da1 ¼ Dx0 þ Dy0 þ þ Dz
@x0 @y0 @z00 0
@d1 @d1 @d1 0
Dd1 ¼ Dx0 þ Dy0 þ þ Dz
@x0 @y0 @z00 0
@a2 @a2 @a2 0
@x0 @y0 @z00 0
@d2 @d2 @d2 0
@x0 @y0 @z00 0
@a3 @a3 @a3 0
@x0 @y0 @z00 0
@d3 @d3 @d3 0
@x0 @y0 @z00 0
When we use more measurements than those strictly necessary, that is, when a
set of n measurements {a1, d1, a2, d2, …, an, dn} is available, where n 3, then
2n independent equations can be written as follows
@a1 @a1 @a1

Da1 ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
@d1 @d1 @d1 0
Dd1 ¼ Dx0 þ Dy0 þ þ 0 Dz0
@x0 @y0 @z0
@a2 @a2 @a2 0
@x0 @y0 @z0
..
.
@an @an @an
Dan ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
@dn @dn @dn
Ddn ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
Thus, the problem reduces to obtaining suitable numerical values of the partial
derivatives, because these values, introduced into the preceding equations, make it
possible to obtain the values of the corrections Dx0, Dy0, …, Dz′0 which best fit the
available measurements a1, d1, a2, d2, …, an, dn.
To this end, Escobal [22] suggests the following procedure. Let e be any one of
the six variables x0, y0, z0, x′0, y′0, and z′0. Let De be some small change introduced
in that variable. The partial derivatives of ai and di (i = 1, 2, …, n) with respect to e
can be approximated by means of the following 2n expressions

@ai ai x0 ; . . .; 0 þ D; . . .; z00 ai x0 ; . . .; 0 ; . . .; z00

@ D

@di di x0 ; . . .; 0 þ D; . . .; z00 di x0 ; . . .; 0 ; . . .; z00

@ D
where e0 is the value of e at the epoch chosen (t0). Each of the variables is incre-
mented in turn, while the others maintain their original values. By so doing, the two
equations written above approximate the needed 6n partial derivatives. Usually an
increment De equal to a few units per cent (Bate et al. [5], suggest 1 or 2%) suffices
to provide a satisfactory approximation of the partial derivatives ∂ai/∂e and ∂di/∂e
(i = 1, 2, …, n). Now the values of these partial derivatives, obtained as has been
shown above, are introduced into
@a1 @a1 @a1

@x0 @y0 @z0
@d1 @d1 @d1
Dd1 ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
@a2 @a2 @a2 0
@x0 @y0 @z0
..
.
@an @an @an
Dan ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
@dn @dn @dn
Ddn ¼ Dx0 þ Dy0 þ þ 0 Dz00
@x0 @y0 @z0
and these equations are solved for Dx0, Dy0, …, Dz′0. Now the variations Dx0, Dy0,
…, Dz′0 are added to the respective quantities x0, y0, …, z′0. This makes it possible
to compute a new couple of position (r) and velocity (r′) vectors
2 3 2 3
x0 þ Dx0 x00 þ Dx00
r0 4 y0 þ Dy0 5 r00 4 y00 þ Dy00 5
z0 þ Dz0 z00 þ Dz00
2.8 Improvement in Orbit Determination by Differential Correction 225
at the same epoch (t0). Hence, new values of right ascension and declination are
computed and compared with the observed values. The differences between the two
sets (computed quantities minus observed quantities) are the residuals. If these
residuals are greater than a chosen tolerance, then the refinement process is repe-
ated; otherwise it is stopped.
2.9 Improvement in Orbit Determination by Weighted

Least-Squares Estimation
The preceding paragraph has shown how to determine the components of the
position and velocity vectors of an observed object so that the residuals should be as
small as possible. In practice, since different measurements have different units and
degrees of reliability, then a weighting factor is applied to each residual, and
consequently the quantity to be minimised is the square of the weighted residuals.
In addition, as shown in Sect. 2.1, the force model used to compute the forces acting
upon the observed object at a time t can only provide approximate values of the true
forces, because the physical constants are known approximately and also because
the mathematical model used to compute the forces can never be exact.
The mathematical formulation of this general problem will be shown below
following Montenbruck and Gill [53]. Let us consider the m-dimensional column
vector x, whose components x1, x2, …, xm are the time-dependent components of
the position (r) and velocity (r′) vectors of the observed object, and the components
of two vectors p and q, which contain the free parameters characterising the force
and measurement model, as follows
2 3
r
6 r0 7
x6
4p5
7
Let
x0 ¼ f ðt; xÞ
x0 ¼ xðt0 Þ
be the differential equation and the initial condition which describe the evolution in
time of the augmented state vector x. Let
2 3
z1
6 z2 7
6 7
z 6 .. 7
4 . 5
zn
be an n-dimensional column vector, whose components z1, z2, …, zn are the

observations at times t1, t2, …, tn. These observations can be expressed as follows
zi ðti Þ ¼ gi ðti ; xi Þ þ ei ¼ hi ðti ; x0 Þ þ ei
(with i = 1, 2, …, n) or in vector terms
z ¼ hðx0 Þ þ e
where gi indicates the computed value of the ith observation as a function of time ti
and the instantaneous state vector xi, and hi indicates the same value as a function of
the state vector x0 at the reference epoch t0. The presence of the quantities ei is due
to the difference between computed and actual values because of measurement
errors, which are assumed to be randomly distributed with zero mean value. Using
the least-squares method, we search the state vector xlsq 0 corresponding to the
minimum value of the loss function
Jðx0 Þ ¼ q q qT q ¼ ½z hðx0 ÞT ½z hðx0 Þ
that is, corresponding to the minimum value of the squared sums of the residuals q1,
q2, …, qn, for a given set of observations z1, z2, …, zn.
The expression of the loss function given above holds for observations of equal
type and quality. This assumption will be removed in the following section of this
paragraph. In addition, the number (n) of observations is assumed to be greater than
or equal to the number (m) of unknowns.
The vector-valued function h(x0) appearing in the equation z = h(x0) + e is a
nonlinear function. Since an approximate value (xappr
0 ) of the actual state vector (x0)
at epoch is known, the problem can be simplified by linearising the function h.
Let xref appr
0 be a reference state vector, which is initially set equal to x0 . The
residual vector q = z − h(x0) is expressed approximately as follows

@h
q ¼ z hðx0 Þ z h xref x0 xref ¼ Dz HDx0
0
@x0 ref
0
where Dx0 = x0 − xref 0 is the difference between x0 and the reference state vector,
Dz = z − h(xref0 ) is the difference between the actual observations and those
resulting from the computed reference orbit, and the Jacobian H = (∂h/∂x0)ref is the
matrix of the partial derivatives of the computed values with respect to the state
vector (here the subscript ref indicates that the partial derivatives are to be evaluated
at x0 = xref
0 ) at the reference epoch t0. Thus, by means of
q ¼ Dz H D x0
we obtain the residual vector q after applying a correction Dx0 to the reference state
vector and recomputing the observations.
2.9 Improvement in Orbit Determination by Weighted Least-Squares Estimation 227
This makes it possible to reduce the original nonlinear least-squares problem of

finding the minimum value of the loss function
Jðx0 Þ ¼ ½z hðx0 ÞT ½z hðx0 Þ
to the simpler linear least-squares problem of searching Dxlsq

0 corresponding to the
minimum value of the loss function
J ðx0 Þ ¼ ½Dz H Dx0 T ½Dz H Dx0
If the Jacobian H has full rank m (the rank of a matrix A is the maximum number
of linearly independent rows or columns of A), that is, if the columns of H are
linearly independent, then the minimum value of the loss function is attained when
the partial derivatives of J(x0) with respect to Dx0 vanish, that is, when
n o
@ ½Dz HDx0 T ½Dz HDx0
¼0
@Dx0
where the partial derivatives are to be evaluated at Dx0 = Dxlsq

0 .
The partial derivatives of the scalar product q q qTq with respect to Dx0 can
be computed by means of the following identity

@ ð aT bÞ @b @a
¼ aT þ bT
@c @c @c
Thus, the general solution of the linear least-squares problem may be obtained
by solving the following system of m algebraic equations
T
HT H Dxlsq
0 ¼ H Dz
The solution is
T 1 T
Dxlsq
0 ¼ H H H Dz
where HTH is an m m symmetric matrix. In order to compute the solution,

Montenbruck and Gill [53] recommend the use of standard techniques for positive
definite linear systems of equations. One of such techniques is the Cholesky
decomposition, which is shown below.
Let A be an m m symmetric (A = AT, where AT is the transpose of A),
positive definite (〈Ax, x〉 > 0 for any real vector x other than the zero vector) matrix
with real entries aij, as follows
2 3
a11 a12 ... a1m
6 a21 a22 ... a2m 7
6 7
A 6 .. .. .. 7
4 . . ... . 5
am1 am2 . . . amm
Then A has a unique decomposition of the form A = LLT, where L is a lower

triangular matrix with positive diagonal entries (‘ii > 0), and LT is the transpose of
L. The two matrices L and LT have the following entries
2 3 2 3
‘11 0 ... 0 ‘11 ‘21 ... ‘m1
6 ‘21 ‘22 ... 0 7 6 0 ‘22 ... ‘m2 7
6 7 6 7
L 6 .. .. .. 7 LT 6 .. .. .. 7
4 . . ... . 5 4 . . ... . 5
‘m1 ‘m2 . . . ‘mm 0 0 ... ‘mm
For example, we want to decompose the following 3 3 symmetric matrix

2 3
4 2 6
A4 2 10 9 5
6 9 26
into the product of two matrices L and LT such that

2 3 2 3
‘11 0 0 ‘11 ‘21 ‘31
L 4 ‘21 ‘22 0 5 LT 4 0 ‘22 ‘32 5
‘31 ‘32 ‘33 0 0 ‘33
By applying the rules of matrix multiplication, we have the following six

equations for the six unknowns ‘11, ‘21, ‘22, ‘31, ‘32, and ‘33:
‘211 ¼ 4 hence ‘11 ¼ 2 ‘221 þ ‘222 ¼ 10 hence ‘22 ¼ 3

‘11 ‘21 ¼ 2 hence ‘21 ¼ 1 ‘21 ‘31 þ ‘22 ‘32 ¼ 9 hence ‘32 ¼ 4
‘11 ‘31 ¼ 6 hence ‘31 ¼ 3 ‘231 þ ‘232 þ ‘233 ¼ 26 hence ‘33 ¼ 1
Consequently, the Cholesky decomposition for A is A = LLT, where

2 3 2 3
2 0 0 2 1 3
L4 1 3 05 LT 4 0 3 4 5
3 4 1 0 0 1
For a 3 3 matrix A, the nonzero entries of L can be computed, as a function of

the entries of A, by means of the following sequence of operations:
1 a21 a31
ð1Þ ‘11 ¼ ða11 Þ2 ð2Þ ‘21 ¼ ð3Þ ‘31 ¼
‘11 ‘11
1 a32 ‘21 ‘31 1
ð4Þ ‘22 ¼ a22 ‘221 2 ð5Þ ‘32 ¼ ð6Þ ‘33 ¼ a33 ‘231 ‘232 2
‘22
In the general case, for the Cholesky decomposition of an m m symmetric,

positive definite matrix A, the following procedure can be used.
For i = 1, 2, …, m and j = i + 1, i + 2, …, m:
!12 !
X
i1
1 X
i1
‘ii ¼ aii ‘2ik ‘ji ¼ aji ‘jk ‘ik
k¼1
‘ii k¼1
Since A is symmetric and positive definite, then the expression under square root
is always positive and the entries ‘ij of L are all real. Therefore, to solve a system
Ax ¼ b
of linear algebraic equations such that A is a square, symmetric and positive definite
matrix, we put the preceding expression in the form
LLT x ¼ b
This done, the equations
Lz ¼ b
are solved for z by forward substitution; then the equations
LT x ¼ z
are solved for x by backward substitution.

Another form of the Cholesky decomposition, which eliminates the need to
compute square roots, is the following
A ¼ LDLT
where A is an m m symmetric, positive definite matrix with real entries aij, D is a

diagonal matrix with all positive nonzero entries djj, and L and LT are, respectively,
a unit lower triangular and a unit upper triangular matrix.
It is to be noted that the matrix L* {‘*ij} of the decomposition A = L*L*T
and the matrix L {‘ij} of the decomposition A = LDLT are not the same matrix.
For example, for m = 3, we have

2 3 2 32 32 3
a11 a12 a13 1 0 0 d11 0 0 1 ‘21 ‘31
4 a12 a22 a23 5 ¼ 4 ‘21 1 0 54 0 d22 0 54 0 1 ‘32 5
a13 a23 a33 ‘31 ‘32 1 0 0 d33 0 0 1
The entries djj of D and the entries ‘ij of L (with i > j) result from
!
X
j1
1 X j1
djj ¼ ajj ‘2jk dkk ‘ij ¼ aij ‘ik ‘jk dkk
k¼1
djj k¼1
For example, we want to decompose the same 3 3 symmetric positive definite

matrix as that considered previously, that is,
2 3
4 2 6
A4 2 10 9 5
6 9 26
into the product LDLT. There results
d11 ¼ a11 ¼ 4
‘21 ¼ a21 =d11 ¼ 2=4 ¼ 1=2
d22 ¼ a22 ‘221 d11 ¼ 10 ð1=2Þ2 4 ¼ 9
‘31 ¼ a31 =d11 ¼ 6=4 ¼ 1:5
‘32 ¼ ð1=d22 Þða32 ‘31 ‘21 d11 Þ ¼ 1=9 ½9 ð3=2Þ 1=2 4Þ ¼ 1:333
d33 ¼ a33 ‘231 d11 ‘232 d22 ¼ 26 ð6=4Þ2 4 ð4=3Þ2 9 ¼ 1
Therefore, the LDLT Cholesky decomposition of the given matrix A is

2 3 2 32 32 3
4 2 6 1 0 0 4 0 0 1 0:5 1:5
4 2 10 9 5 ¼ 4 0:5 1 0 54 0 9 0 54 0 1 1:333 5
6 9 26 1:5 1:333 1 0 0 1 0 0 1
The L*L*T and LDLT Cholesky decompositions of the same matrix A are
related to each other as shown below
1 1 1
1 T 1
1 T
A ¼ LDLT ¼ LD2 D2 LT ¼ LD2 D2 LT ¼ LD2 LD2
where D½ is also a diagonal matrix, whose nonzero entries are the square roots of
the corresponding entries of D, that is, D½ = diag(d½ ½ ½
11, d22, …, dmm). Hence,
A ¼ LDLT ¼ L LT
where L* = L D½.
Now, let us come back to the linear least-squares problem. Since h(x0) is actually
a nonlinear function, the value Dxlsq
0 computed as shown above (and consequently
xlsq ref lsq

0 = x0 + Dx0 ) is only an approximate solution of the orbit determination
problem. A better solution than this can be obtained by replacing xref0 with the value
of xlsq
0 computed previously and performing a new iteration. Let the superscripts
[i] and [i + 1] denote the value of xlsq
0 resulting from, respectively, the ith and the
th
(i + 1) iteration. The iterative process of refinement can be expressed as follows
1 h i
½i þ 1 ½i ½i
x0 ¼ x0 þ H½iT H½i H½iT z h x0
which formula has xappr0 as its starting point, that is, x[0] appr
0 = x0 , and continues till
the difference x0 − x0 is greater in magnitude than a desired tolerance.
[i+1] [i]
For the best convergence of the iterative process, the Jacobian H[i] must be
computed at each iteration; however, Montenbruck and Gill [53] observe that H[i]
can be replaced by a constant H[0]. This replacement implies a higher number of
iterations, but results often in less computational work than would be required if
H[i] were computed at each step. The method shown above does not take into
account the different errors by which the measurements may be affected. However,
this method can be made general by weighting each observation zi with the inverse
of its mean measurement error ri (where i = 1, 2, …, n).
By so doing, each residual qi is replaced by the corresponding normalised
residual q*i defined as follows
qi ½z hðx0 Þi
qi ¼ ¼
ri ri
The mean measurement error ri includes both the random and systematic errors.
A typical example of the latter type of error is provided by the refraction of light
rays due to the Earth atmosphere. By using the normalised residuals q*i instead of
−1
the ordinary residuals qi, the expression Dxlsq T
0 = (H H) (HT Dz) shown above
remains formally identical, the only difference being that in the present case
(weighted observations) the Jacobian H and the difference vector Dz must be
replaced by their respective counterparts H* and Dz*, as follows
T 1 T
Dxlsq
0 ¼ H H H Dz
where H* = R H and Dz* = R Dz; R is the n n diagonal matrix filled with zeros,
except the elements placed along its main diagonal, which are r−1 −1
1 , …, rn .
−1 −1
In other words, R diag(r1 , …, rn ) or, which is the same,
2 3
1=r1 0 ... 0
6 0 1=r2 ... 0 7
6 7
R6 . .. .. 7
4 .. . ... . 5
0 0 ... 1=rn
The solution given above

T 1 T
Dxlsq
0 ¼ H H H Dz
of the problem of weighted least-squares estimation may also be expressed in terms

of the original Jacobian matrix H and the original difference vector Dz (instead of
H* = R H and Dz* = R Dz) as follows
T 1 T
Dxlsq
0 ¼ H WH H W Dz
where W is the weighting matrix defined as W R2 diag(r−2 −2

1 , …, rn ), or,
which is the same,
2 3
1=r21 0 ... 0
6 0 1=r22 ... 0 7
6 7
W R2 6 . .. .. 7
4 .. . ... . 5
0 0 ... 1=r2n
The preceding definition of W as a diagonal matrix holds in case of uncorrelated

errors of measurement. Should such errors be correlated, then W would become a
non-diagonal matrix. In order to understand the correlation of errors, it is necessary
to have some concepts of probability theory, which are given below.
As has been shown above, if x0 and e designate, respectively, the actual (aug-
mented) state vector at epoch and the vector containing the measurement errors,
then the observation vector z results from
z ¼ hðx0 Þ þ e
The preceding expression, after linearisation, becomes

Dz ¼ H x0 xref
0 þe
where xref
0 is a reference state vector sufficiently close to x0. The solution of this
least-squares problem has been shown to be
T 1 T
xlsq
0 ¼ x0 þ H WH
ref
H WDz
1 T
¼ x0 þ HT WH H We
The preceding expression shows that, when measurement errors are committed,
the computed state vector xlsq
0 differs from the actual state vector x0.
In the event of systematic errors being negligible, the components of e are only
random errors. Let X be a discrete random variable which takes values in S = {x1,
x2, …, xn}. The expected value (or mean) of X is defined as follows
X
n
EðXÞ ¼ xi pðxi Þ
i¼1

where p(xi) is the probability of xi. The mean of X is also denoted by X or by lX.
The preceding definition holds if the sum x1 p(x1) + x2 p(x2) + + xn p(xn)
converges to a finite value. Otherwise, the expected value of X is undefined.
For example, when a die is cast, the probability of each of the six possible events
xi (with i = 1, 2, …, 6) is p(xi) = 1/6. Thus, according to the preceding definition,
the expected value is E(X) = (1 + 2 + 3 + 4 + 5 + 6) (1/6) = 21/6 = 7/2 = 3.5.
Let X be a continuous random variable in S = [−∞, ∞] with probability density
function u(x). In this case, the expected value of X is defined as follows
Z1
EðXÞ ¼ x uðxÞ dx
1
Here, too, the value of the integral must be finite, in order for the expected value
of X to be defined. Let X and E(X) be, respectively, a random variable and its
expected value. The variance of X is defined as follows
n o
VarðXÞ ¼ E ½X EðXÞ2
In other words, the variance of X is the expected squared distance of X from its
mean E(X). Let X be a discrete random variable which takes values in S = {x1, x2,
…, xn}. Then the variance of X is defined as follows
" #2
X
n X
n
Var ð X Þ ¼ xi x j p x j pð x i Þ
i¼1 j¼1
In the example of the die proposed above, p(xi) = 1/6 and E(X) = 7/2.
Thus, according to the preceding definition, the variance of X is
"
7 2 7 2 7 2 7 2 7 2
Var ð X Þ ¼ 1 þ 2 þ 3 þ 4 þ 5
2 2 2 2 2
#
7 2 1 25 9 1 1 9 25 1 35
þ 6 ¼ þ þ þ þ þ ¼
2 6 4 4 4 4 4 4 6 12
Likewise, let X be a continuous random variable in S = [−∞, ∞] with proba-

bility density function u(x). Then the variance of X is
2 32
Z1 Z1
Var ð X Þ ¼ 4x nu ðnÞ dn5 u ðxÞ dx
1 1
In case of the expected value E(X) of a random variable X being zero, that is, in
case of E(X) = 0, then the variance of X is
n o h i h i
VarðXÞ ¼ E ½X EðXÞ2 ¼ E ðX 0Þ2 ¼ E ðXÞ2
The variance of a random variable X is often denoted by r2X, that is,

Var(X) r2X, and its standard deviation [Var(X)]½ by rX.
Let X and Y be two independent random variables. Let E(X) and E(Y) be their
respective expected values. In this case, there results
VarðX þ YÞ ¼ VarðXÞ þ VarðYÞ
Instead, in case of X and Y being dependent random variables, then there results
VarðX þ YÞ ¼ VarðXÞ þ Var ðY Þ þ 2 E f½X EðXÞ½Y EðYÞg
where the term E{[X − E(X)][Y − E(Y)]} is called the covariance of X and Y, that
is,
Cov ðX; YÞ ¼ E f½X EðXÞ½Y EðYÞg
The covariance of two random variables X and Y is the average value of the
deviation of X from its mean E(X) and the deviation of Y from its mean E(Y). The
zero value of the covariance Cov(X, Y) of two random variables X and Y does not
imply their independence; it implies only the linear independence of X and Y.
The correlation coefficient qXY Corr(X, Y) of two random variables X and Y is
defined as follows
CovðX; Y Þ
Corr ðX; Y Þ ¼ 1
½VarðXÞVarðYÞ2
Corr(X, Y) is a dimensionless quantity having the following properties:
1 Corr ðX; Y Þ 1
Corr ðaX þ b; cY þ d Þ ¼ Corr ðX; Y Þ
where a, b, c, and d are any real constants. This implies that Corr(X, Y) is equal to
unity only when a relation of a linear dependence exists between X and Y.
Let us see now the practical application of the concepts given above. If the error
vector e has only random errors as its components, then the expected values of,
respectively, e and eeT are such that
EðeÞ ¼ 0
EðeeT Þ ¼ diag ðr21 ; . . .; r2n Þ
In other words, the expected value E(ei) of each component ei (i = 1, 2, …, n) of

the error vector e is equal to zero; all these components are uncorrelated, in other
words ei ej = 0 for i 6¼ j; and the standard deviation ri of ei is equal to [E(e2i )]½.
The least-squares solution xlsq
0 of the orbit determination problem is a random
variable, because it depends on e, which is a random variable by the hypothesis
made above. Then, E(e) = 0 implies that the expected value of xlsq 0 is equal to the
actual value of x0, because
h 1 T i 1 T
E xlsq
0 ¼ E x0 þ HT W H H W e ¼ x0 þ H T W H H W EðeÞ ¼ x0
In addition, since by definition Cov(X, Y) = E{[X − E(X)][Y − E(Y)]}, then

h ih iT
Cov xlsq lsq
0 ; x0 ¼E xlsq
0 E xlsq
0 xlsq
0 E xlsq
0
h ih i T
¼E xlsq
0 x 0 xlsq
0 x0
Substituting
1 T
x0 þ HT W H H We
for xlsq
0 into the preceding expression yields
h ih iT
T 1 T T 1 T
Cov xlsq lsq
0 ; x0 ¼ E x0 þ H W H H W e x0 x0 þ H W H H W e x0
h i h i
T 1 T T 1 T T
¼E H WH H We H WH H We
1 T T 1
¼ HT W H H W E ee ðW HÞ HT W H
The expression written above can be simplified, if the weighting matrix W is

chosen in accordance with the standard deviation ri (where i = 1, 2, …, n) of the
measurement. In this case, W = diag(r−2 −2
1 , …, rn ) is the inverse of E(ee ), and the
T
preceding expression reduces to

1
Cov xlsq lsq
0 ; x0 ¼ HT W H
Those elements of the covariance matrix which are placed along its main
diagonal are the standard deviation
h i12
r xlsq
0i ¼ Cov xlsq lsq
0i ; x0i
(i = 1, 2, …, n) of the components of the xlsq 0 vector. Likewise, the off-diagonal

elements of the covariance matrix provide a measure of the correlation existing
between errors of individual components.
The expected value and the covariance of xlsq
0 define the distribution of values of
lsq
x0 which would result in an experiment of repeated orbit determinations for the
same trajectory, if the measurement errors were only of the random type. If these
errors have a normal distribution, then there is a probability of 67% that xlsq 0
(resulting from the actual measurements) deviates from x0 by less than 1r, and a
probability of 99.7% that xlsq
0 deviates from x0 by less than 3r.
In the presence of systematic errors e*, there is a further deviation
T 1 T
d xlsq
0 ¼ H WH H W e
of xlsq
0 from x0. Montenbruck and Gill [53] point out that the measurement standard
deviation r(e) must be known in order to construct the weighting matrix W.
The analysis performed so far is based on the assumption of a Gaussian distri-
bution of errors in the observations. However, this analysis (based only on data
noise errors) does not take into account the effect of model errors. It is therefore
necessary to take into account the effect due to systematic errors.
To this end, the measurement equation can be rewritten as follows
z ¼ hð x 0 ; c Þ þ e
where z is the n-dimensional vector containing the observations at times t1, t2, …,
tn; x0 is the vector of the estimated parameters; c is the vector of the so-called
consider parameters (which affect the estimated parameters); h is a vector-valued
function of x0 and c; and e is the vector containing the measurement noise errors.
The vector c contains the force and measurement parameters which are uncertain
but are not modified as a result of the least-squares estimation process. The
expected value of these parameters can be assumed, without loss of generality,
equal to zero. Previously both the estimated and the consider parameters have been
taken together and the expression
z ¼ hðx0 Þ þ e
of the observation vector z has been linearised around the reference state vector xref
0
to obtain

Dz ¼ H x0 xref0 þe
Just in the same way, now the expression
z ¼ hð x 0 ; c Þ þ e
is linearised around xref

0 to obtain

Dz ¼ Hx x0 xref
0 þ Hc c þ e
where Hx is the Jacobian matrix containing the partial derivatives of the

vector-valued function h with respect to x0, and Hc is the Jacobian matrix con-
taining the partial derivatives of h with respect to c. Now the expression of the
least-squares solution becomes
T 1 T
xlsq
0 ¼ x0 þ H x W H x H x WðH c c þ e Þ
This solution differs from the true value x0 of the estimation parameters by a
quantity which depends on the consider parameters (components of c) and the
measurement noise (components of e). The consider parameters are assumed to be
random variables with zero mean and covariance matrix C, which are uncorrelated
with the measurement noise, so that the equality E(ceT) = 0 holds. In this case, the
expected value of the least-squares solution
1
E xlsq
0 ¼ x0 þ HTx W Hx HTx W½Hc EðcÞ þ EðeÞ ¼ x0
is equal to the true state vector x0.

The consider covariance matrix Pc is larger than the noise-only covariance
matrix P = (HTx W Hx)−1, which is also called formal or computed covariance
matrix. The consider covariance matrix has the following expression
T
Pc ¼ H HTx W Hc C HTc þ EðeeT Þ P HTx W
T
¼ P þ ðP HTx WÞðHc C HTc Þ P HTx W
where the weighting matrix W has been taken as the inverse of the measurement
covariance matrix.
As has been shown above, an approximate value (xappr
0 ) of the actual state vector
(x0) at epoch must be known to start the least-squares orbit determination process.
Some information on the accuracy of this value is often available. We want to
incorporate the a priori covariance matrix Papr
0 into the least-squares estimation. To
his end, the loss function shown above
J ðx0 Þ ¼ q q qT q ¼ ðDz H Dx0 ÞT ðDz H Dx0 Þ
is represented in another way. Remembering that

T 1 T
Dxlsq
0 ¼ H H H Dz
the loss function may be written as follows

T
J ðx0 Þ ¼ Dx0 Dxlsq0 ðHT
HÞ Dx 0 Dx lsq
0 þ Dz T
Dz Dx lsq T T
0 H H Dx lsq
0
T
¼ x0 xlsq
0 P0 1 x0 xlsq
0 þ constant
Since this expression of the loss function J(x0) is a quadratic form of (x0 − xlsq
0 )
defined by the inverse covariance matrix P−1 0 = (H H) of (x0 − x0 ), then the loss
T lsq
function minimum and the covariance matrix provide the same information on the
least-squares estimation as that which comes from the measurement vector Dz and
the Jacobian matrix H. Thus, an a priori estimate
xapr apr
0 ¼ x0 þ Dx0
ref
of the state vector x0 may come from a modified loss function

T
J ¼ x0 xapr
0 K x0 xapr
0 þ qT q
−1
where K = (Papr 0 ) , called information matrix, is used to attribute a contribution to
the loss function to each deviation from xapr 0 , and q is the vector containing the
residuals. The information matrix K is always positive semi-definite, because it is
the inverse of the covariance matrix. By the way, let A be a matrix with real entries
aij. A is said to be positive semi-definite if, for any vector x with real components xi,
the dot product 〈Ax, x〉 of Ax and x is non-negative, that is, if
hA x; xi 0
Consequently, K can be factored to form a product K = STS, which makes it

possible to determine the minimum of the loss function J. If J is written as follows
J ¼ ðDx0 Dxapr T apr T

0 Þ KðDx0 Dx0 Þ þ ðDz H Dx0 Þ ðDz H Dx0 Þ ¼ A A
T

SDxapr S
A¼ 0 Dx0
Dz H
then the information matrix K can be considered as the result of additional

observations. Consequently, the minimum of the combined loss function results,
after simplification, from
1
Dx0 lsq ¼ K þ HT H K Dxapr
0 þ H Dz
T
In case of weighted observations, the preceding expression becomes

1
Dxlsq
0 ¼ KþH WH
T
K Dxapr
0 þ H W Dz
T
where W is the weighting matrix. In order for the preceding expression to be

computable, the sum K + HTWH must have a nonzero determinant, without the
need for K or HTWH to be non-singular (by the way, an n n matrix A is said to
be non-singular or invertible when there exists an n n matrix B A−1 such that
AB = BA = I, where I is the n n identity matrix). However, the non-singularity
of the information matrix K is sufficient to ensure the resolvability of the equations
independently of HTWH. In order to take advantage of this fact, a singularity in
least-squares problems can be avoided by giving a small a priori weight to each
estimation parameter and adding the corresponding diagonal matrix K to the normal
equations matrix.
−1
The expected value Dxlsq0 = (K + H WH)
T
(K Dxapr T
0 + H W Dz) of the esti-
mated state is equal to the actual state x0, if the a priori information xapr
0 is also a
random variable having x0 as its mean value.
The covariance matrix P0 of the estimate is related to the a priori covariance and
the measurement information matrix by the following expression
ðP0 Þ1 ¼ ðPapr

0 Þ
1
þ ðHT WHÞ
2.10 Numerical Solution of the Least-Squares Estimation

Problem
The purpose of the present paragraph is to provide the reader of this book with the
basic concepts which are necessary to solve numerically the least-squares problem
described in the preceding paragraph. Those among the readers who are fully
conversant with such concepts can skip this paragraph and go directly to the next.
For the sake of generality, the system of m normal equations of Sect. 2.9, that is,
(HT H)Dxlsq T
0 = (H Dz), is written here in the following form

AT A x ¼ AT b
where A, x, and b take the places of, respectively, H, Dxlsq

0 , and Dz.
The numerical procedures to be described below are based on the decomposition
(also called factorisation) of a given non-singular n m matrix A (with n m)
into an orthogonal n n matrix Q and an upper triangular m m matrix R. As the
bottom (n − m) rows of an n m upper triangular matrix contain only zeroes, as

will be shown below, it is customary to write

Rmm
Anm ¼ Qnn
0ðnmÞm
where 0 is an (n − m) m partition, containing only zeroes, of the given matrix A.

According to the definition given by Olver and Shakiban [60], an orthogonal (or
orthonormal) matrix Q is a square matrix which satisfies the condition QTQ = I,
where QT is the transpose of the given matrix and I is the identity matrix. Since the
inverse matrix Q−1 must satisfy the condition Q−1 Q = I, then a square matrix Q is
orthogonal if and only if its transpose QT is equal to its inverse Q−1.
In other words, a square matrix is orthogonal if and only if its columns form an
orthonormal basis with respect to the Euclidean scalar (or inner or dot) product on
an n-dimensional Euclidean space. This is because, if q1, q2, …, qn are the columns
of Q, then qT1 , qT2 , …, qTn will be the rows of the transpose matrix QT. Now, the (i, j)
entry of the product QTQ results from the product between the ith row of QT and the
jth column of Q. Since Q is, by hypothesis, an orthogonal matrix (such that
QTQ = I), then the following relations must hold

1 ði ¼ jÞ
qi qj ¼ qTi qj ¼
0 ði 6¼ jÞ
which are just the conditions for q1, q2, …, qn to form an orthonormal basis.
Guillemin [31] gives the following numerical example of an orthogonal matrix
of order three
2 3
0:5 0:5 0:707
Q 4 0:707 0:707 0 5
0:5 0:5 0:707
where 0.707 approximates ½(2)½. As is easy to verify, the three columns

2 3 2 3 2 3
0:5 0:5 0:707
q1 4 0:707 5 q2 4 0:707 5 q3 4 0 5
0:5 0:5 0:707
of the 3 3 matrix Q indicated above satisfy the conditions qi qj = 1 (for i = j)

and qi qj = 0 (for i 6¼ j). An orthogonal matrix Q has determinant det(Q) = ± 1.
This is because Q, as an orthogonal matrix, must satisfy the condition QTQ = I;
taking the determinant of this equality yields

1 ¼ det ðIÞ ¼ detðQT QÞ ¼ det QT det ðQÞ ¼ ½det ðQÞ2
If the determinant of an n n orthogonal matrix is equal to +1, then that matrix

is called proper and the corresponding orthonormal basis is a right-handed basis on
2.10 Numerical Solution of the Least-Squares Estimation Problem 241
an n-dimensional Euclidean space; whereas, if the determinant of an orthogonal

matrix is equal to −1, then that matrix is called improper and the corresponding
orthonormal basis is a left-handed basis on the same space.
The product of two orthogonal matrices is also an orthogonal matrix. This is
because, if Q1 and Q2 are two orthogonal matrices, that is, two matrices such that
QT1 Q1 ¼ I
QT2 Q2 ¼ I
then there results (Q1 Q2)T (Q1 Q2) = QT2 QT1 Q1 Q2 = QT2 Q2 = I. Thus, the
product matrix Q1 Q2 is also orthogonal.
An m m matrix R is said to be upper triangular, if its entries rij below the main
diagonal are zero (rij 6¼ 0 for i j; rij = 0 for i > j), as shown below:
2 3
r11 r12 ... r1m
6 0 r22 ... r2m 7
6 7
R 6 .. .. .. 7
4 . . ... . 5
0 0 . . . rmm
The QR decomposition applied to a given non-singular matrix A makes the

product ATA less sensitive to small errors affecting the values of the entries aij of A.
The QR decomposition is unique, if all the diagonal elements rii (i = 1, 2, …, m) of
R are required to be real and positive [83].
There are several methods for computing the QR decomposition. Among them
are the Gram–Schmidt procedure, the Householder transformations, and the Givens
rotations. These three methods will be described below. The interested reader can
find more on the matter in several books or articles, for example in Refs. [53, 29].
The numerically stable Gram–Schmidt procedure is a method to make a set of
vectors orthogonal in an inner product space.
An inner product space is a vector space with an operation, which associates
each pair v and w of vectors of the space with a scalar quantity known as the inner
product 〈v, w〉 of the considered pair of vectors. A vector space V is the set of all
real (column) vectors v with n components v1, v2, …, vn, this set being closed under
the two operations of addition (if v and w are two vectors of V, then v + w is also a
vector of V) and multiplication by a scalar (if v is a vector of V and s is a real scalar,
then sv is also a vector of V). Two vectors v and w of a vector space V are said to be
orthogonal, if their inner product 〈v, w〉 is zero. Vectors v1, v2, …, vn of an n-
dimensional vector space V are linearly independent, if the equality c1v1 +
c2v2 + + cnvn = 0 can only hold if the coefficients c1, c2, …, cn are all equal to
zero. Vectors v1, v2, …, vn of an n-dimensional vector space V are a basis, if any
vector v of V can be expressed as a linear combination v = c1v1 + c2v2 + + cnvn
with a unique choice of the coefficients c1, c2, …, cn. For an n-dimensional vector
space, any nonzero linearly independent vectors form a basis. Let V and W be two
sets of vectors. If each vector v contained in V can be expressed as a linear
combination of the vectors contained in W, then W is said to be the basis set (or the
generating set or the spanning set) for V.
The Gram–Schmidt procedure takes a finite set {a1, a2, …, am} of linearly
independent vectors and generates an orthogonal set {q1, q2, …, qm} which spans
the same subspace as the previous set. Let
2 3
a11 a12 ... a1m
6 0 a22 ... a2m 7
6 7
A 6 .. .. .. 7
4 . . ... . 5
an1 an2 ... anm
be a given n m matrix with m linearly independent columns. The same matrix

can also be indicated briefly A [a1 a2 … am], where a1, a2, …, am are the column
vectors of A. We want to construct an orthogonal n n matrix Q {q1, q2, …,
qn} and an upper triangular m m matrix R {r1, r2, …, rm} such that

Rmm
Anm ¼ Qnn
0ðnmÞm
For example, let A be a two-column matrix A [a1 a2] and Q be a two-column

matrix Q [q1 q2]. Since R must be a triangular matrix, then

r11 r12
A ½ a1 a2 ¼ ½ q1 q2 ¼ QR
0 r22
In this case, the two column vectors q1 and q2 of Q and the three nonzero entries
r11, r12, and r22 of R are the quantities to be determined as a function of the two
column vectors a1 and a2 of the given matrix A.
Likewise, when A is a three-column matrix A [a1 a2 a3], then
2 3
r11 r12 r13
A ½ a1 a2 a3 ¼ ½ q1 q2 q3 4 0 r22 r23 5 ¼ QR
0 0 r33
In this case, the three column vectors q1, q2, and q3 of Q and the six nonzero
entries r11, r12, r13, r22, r23, and r33 of R are the quantities to be determined as a
function of the three column vectors a1, a2, and a3 of the given matrix A.
The numerically stable Gram–Schmidt procedure is described below by means
of an example, due to Wikipedia [82]. Let us consider the following 3 3 matrix
2 3
12 51 4
A4 6 167 68 5
4 24 41
The three column vectors a1, a2, and a3 of the given matrix A will, by degrees,
be replaced by the three column vectors q1, q2, and q3 of an orthogonal matrix Q.
At the same time, the six nonzero entries r11, r12, r13, r22, r23, and r33 of an upper
triangular matrix R will be computed so that A = QR.
To this end, we first compute the norm (that is, the Euclidean length) of the first
(or leftmost) column vector a1 of A. The first entry r11 of R is set equal to the norm
||a1||. This yields
1 h i12
r11 ¼ ka1 k ¼ a211 þ a221 þ a231 2 ¼ 122 þ 62 þ ð4Þ2 ¼ 14
Now, we normalise a1 by dividing all its components by its norm ||a1||. This
yields the first column q1 of Q, as follows
a11 12 6
q11 ¼ ¼ ¼
ka1 k 14 7
a21 6 3
q21 ¼ ¼ ¼
ka1 k 14 7
a31 4 2
q31 ¼ ¼ ¼
ka1 k 14 7
At this stage of the procedure, the given matrix A becomes

2 3
6=7 51 4
4 3=7 167 68 5
2=7 24 41
Now, we compute the next entries (r12, r13, …, r1m) of R by taking the scalar
products of the first column of the matrix shown above with each of the other
columns of the same matrix. In the example considered above, this yields
r12 ¼ ð6=7Þ ð51Þ þ ð3=7Þ 167 þ ð2=7Þ 24 ¼ 21

r13 ¼ ð6=7Þ 4 þ ð3=7Þ ð68Þ þ ð2=7Þ ð41Þ ¼ 14
Now, for the purpose of making the second, the third, …, the mth column
orthogonal to the first column, we subtract r12 times the first column from the
second column, r13 times the first column from the third column, …, r1m times the
first column from the mth column. This yields
51 ð6=7Þ 21 ¼ 69 4 ð6=7Þ ð14Þ ¼ 16

167 ð3=7Þ 21 ¼ 158 68 ð3=7Þ ð14Þ ¼ 62
24 ð2=7Þ 21 ¼ 30 41 ð2=7Þ ð14Þ ¼ 45
and the resulting matrix is

2 3
6=7 69 16
4 3=7 158 62 5
2=7 30 45
At this stage of the procedure, the first column vector of the preceding matrix has
been normalised to have unit length; in addition, the second column and the third
column have been made orthogonal to the first.
Now, in order to compute q2, we normalise the second column, by dividing all
its components by its norm. At the same time, the r22 entry of R is set equal to the
norm (or Euclidean length) of the second column. In the example considered above,
this yields
h i12
r22 ¼ ð69Þ2 þ 1582 þ 302 ¼ 14
The resulting matrix is

2 3
6=7 69=175 16
4 3=7 158=175 62 5
2=7 6=35 45
Now r23 is the scalar product of the second column and the third column:

69 158 6
r23 ¼ 16 þ ð62Þ þ ð45Þ ¼ 70
175 175 35
By subtracting r23 times the second column from the third, we obtain
16 ð70Þ ð69=175Þ ¼ 406=35

62 ð70Þ ð158=175Þ ¼ 42=35
45 ð70Þ ð6=35Þ ¼ 33
The resulting matrix is

2 3
6=7 69=175 406=35
4 3=7 158=175 42=35 5
2=7 6=35 33
Now we compute r33 as the norm of the third vector:

" 2 2 #12
406 42 2
r33 ¼ þ þ ð33Þ ¼ 35
35 35
Finally, in order to compute q3, the last column is normalised by dividing all its
components by its norm. This yields

406 1 58
¼
35 35 175

42 1 6
¼
35 35 175
1 33
ð33Þ ¼
35 35
The two matrices Q and R resulting from the procedure described above are then
2 3
6=7 69=175 58=175
Q 4 3=7 158=175 6=175 5
2=7 6=35 33=35
2 3
14 21 14
R4 0 175 70 5
0 0 35
where A = QR. The example shown above has described the numerically stable
Gram–Schmidt procedure for a given 3 3 matrix A. The same procedure can be
extended to any non-singular m m matrix A, by means of the algorithm given
below, which is taken from Olver and Shakiban [60]. This algorithm takes the
entries aij of A and replaces them with the entries qij of Q; at the same time, it
computes the nonzero entries rij of R, which must be stored in a separate matrix.
start
for j = 1 to m
set rjj = (a21j + a22j + … + a2mj)½
if rjj = 0, stop; print “A has linearly dependent columns”
else for i = 1 to m
set aij = aij/rjj
next i
for k = j + 1 to m
set rjk = a1ja1k + a2ja2k + … + amjamk
for i = 1 to m
set aik = aik – aijrjk
next i
next k
next j
end
The QR decomposition based on the Householder transformations is shown

below. Let x and y be any two linearly independent vectors, having the same
Euclidean length ||x|| = ||y||, of an n-dimensional vector space V.
Consider the unit vector u = (y − x)/||y − x|| and the matrix H = I − 2(u uT).
Then H is the reflection matrix such that H x = y. It is to be noted that uuT is an
outer (or tensor) product, which yields a matrix, not a scalar.
For example, let
2
3 2 3
12 14
x4 6 5 y4 0 5
4 0
be two column vectors of a three-dimensional vector space V.

It is easy to verify that the two vectors x and y given above are linearly independent
and have the same Euclidean length ||x|| = ||y|| = 14. In this case, there results
2 3
2
y x ¼ 4 6 5
4
h i12 1
ky xk ¼ 22 þ ð6Þ2 þ 42 ¼ 2ð14Þ2
2 1
3
1=ð14Þ2
yx 6 7
u¼ ¼6 3=ð14Þ2 7
1
ky xk 4 5
1
2=ð14Þ2
2 1
3
1=ð14Þ2 h i
6 7
uuT ¼ 6 7
1 1 1 1
4 3=ð14Þ2 5 1=ð14Þ2 3=ð14Þ2 2=ð14Þ2
1
2 2=ð14Þ
2
3
1=14 3=14 2=14
6 7
¼ 4 3=14 9=14 6=14 5
2=14 6=14 4=14
2 3
1=7 3=7 2=7
6 7
2 uu T
¼ 4 3=7 9=7 6=7 5
2=7 6=7 4=7
2 3 2 3
1 0 0 1=7 3=7 2=7
6 7 6 7
H ¼ I 2 uuT ¼ 4 0 1 0 5 þ 4 3=7 9=7 6=7 5
0 0 1 2=7 6=7 4=7
2 3
6=7 3=7 2=7
6 7
¼ 4 3=7 2=7 6=7 5
2=7 6=7 3=7
As is easy to verify, the reflection matrix H, found as shown above, satisfies the
condition H x = y. In addition, H is an orthogonal matrix (that is, HTH = I) and its
determinant det(H) is equal to −1. Since we know how to compute the reflection
matrix H in an n-dimensional vector space, we can apply these concepts to produce
QR decompositions. To this end, we intend to introduce sub-diagonal zeros into the
given matrix A to be decomposed, in order to gradually change A into an upper
triangular matrix R.
In other words, we multiply the given matrix Anm on the left by a sequence of
reflection matrices H1, H2, …, Hk, so that the product Hk H2 H1 A should be
equal to an upper triangular matrix R. The first step operates on the matrix A itself.
We choose the first (that is, the leftmost) column vector a1 of the given matrix
Anm [a1 a2 … am]. We compute the Euclidean length
1
ka1 k ¼ a211 þ a221 þ þ a2n1 2
of the first column vector a1 and then find a reflection matrix H1 such that the first
column of the product H1 A should be a multiple of the vector [1 0 0 … 0]T. In
other words, we form
2 3
1
607
6 7 v
v ¼ a1 ka1 k 6 07
6.7 u¼ H1 ¼ I 2ðu uT Þ
4 .. 5 kv k
0
By so doing, H1 is a reflection matrix such that the first column of H1 A is a

multiple of the vector [1 0 0 … 0]T. In other words, the first column of H1 A is a
column vector having all zeros in its rows except the first.
Each of the column vectors a1, a2, …, ai, …, am of Anm can be reflected onto a
multiple of [1 0 0 … 0]T in two ways: it can be reflected onto ||ai|| [1 0 0 … 0]T or
onto −||ai|| [1 0 0 … 0]T. The choice of the sign in front of ||ai|| is important, because
when we form the unit vector u = (y − x)/||y − x|| we divide by ||y − x||.
Consequently, a choice which makes this denominator small must be avoided. To
this end, the sign in front of ||ai|| should be the opposite of the sign in front of the
entry placed in the kth row of ||ai||, where aki is the pivot co-ordinate after which all
entries are zero in the final upper triangular form of A. Therefore, it is advisable to
choose
u ¼ sgn ðaki Þkai k½ 1 0 0 . . . 0 T
where the signum function sgn(x) is such that

8
< 1 ðif x\ 0Þ
sgnðxÞ ¼ 0 ðif x ¼ 0Þ
:
1 ðif x [ 0Þ
The second step operates on the matrix A′, which results from cancelling the first
row and the first column of H1A and retaining the rest of H1A. We repeat for A′ the
same operations performed in the first step, and compute the reflection matrix H′2.
Since H′2 is of smaller rank than H1 and we want to operate with H1A (and not with
A′), then we expand A′, by filling in a 1 in the upper left entry of H1A. This means
that the second reflection matrix H2 results from

1 0
H2 ¼
0 H0 2
The third step operates on A′′, which results from cancelling the first row and the
first column of H2A′ and retaining the rest of H2A′. We repeat for A′′ the same
operations performed in the second step, and compute the reflection matrix H′3.
The third reflection matrix H3 results from
2 3
1 0 0
H3 ¼ 4 0 1 0 5
0 0 H0 3
In general, at the pth step, the pth reflection matrix Hp results from

Ip1 0
Hp ¼
0 H0 p
After a number k = min(m − 1, n) of steps, the result of this process will be
Hk H2 H1 A ¼ R
where R is an upper triangular matrix, and each of the reflection matrices (H1, H2,
…, Hk) is an orthogonal matrix. Thus, with
Q ¼ H1 T H2 T HTk
the desired result A = QR will be reached. For example, we consider the matrix
2 3
12 51 4
A4 6 167 68 5
4 24 41
which was decomposed previously by means of the Gram–Schmidt procedure.

We consider the first column vector a1 [12 6 −4]T of the matrix A and
compute the Euclidean length ||a1|| of a1, as follows
h i12
ka1 k ¼ 122 þ 62 þ ð4Þ2 ¼ 14
Then we compute
k a1 k½ 1 0 0 1 T ¼ ½ 14 0 0 T
v ¼ a1 sgnða11 Þka1 k½ 1 0 0 0 T ¼ ½ 2 6 4 T
h i12 1
kvk ¼ ð2Þ2 þ 62 þ ð4Þ2 ¼ 2ð14Þ2
2 1
3
1=ð14Þ2
v 6 1 7
u¼ ¼ 4 3=ð14Þ2 5
kv k 1
2=ð14Þ2
2 1
3
1=ð14Þ2 h i
6 1 7
uuT ¼ 4 3=ð14Þ2 5 1=ð14Þ2 3=ð14Þ2
1 1 1
2=ð14Þ2
1
2 2=ð14Þ
2
3
1=14 3=14 2=14
¼ 4 3=14 9=14 6=14 5
2=14 6=14 4=14
2 3 2 3
1 0 0 1=7 3=7 2=7
T 6 7 6 7
H1 ¼ I 2 uu ¼ 4 0 1 0 5 þ 4 3=7 9=7 6=7 5
0 0 1 2=7 6=7 4=7
2 3
6=7 3=7 2=7
6 7
¼ 4 3=7 2=7 6=7 5
2=7 6=7 3=7
Now we consider the matrix resulting from the product H1A, that is,
2 32 3 2 3
6=7 3=7 2=7 12 51 4 14 21 14
H1 A ¼ 4 3=7 2=7 6=7 54 6 167 68 5 ¼ 4 0 49 14 5
2=7 6=7 3=7 4 24 41 0 168 77
This matrix, due to the nonzero value of its (H1A)32 entry (because 168 6¼ 0), is
not an upper triangular matrix. Therefore, we cancel the first row and the first
column of H1A, and consider the matrix

49 14
A0 ¼
168 77
The Euclidean length of the first column vector a′1 of A′ [a′1 a′2] results from
0 h i1
a ¼ ð49Þ2 þ 1682 2 ¼ 175
1

49 1 224
v¼ 175 ¼
168 0 168
h i12
kvk ¼ ð224Þ2 þ 1682 ¼ 280

v 4=5
u¼ ¼
kvk 3=5

4=5 16=25 12=25
uuT ¼ ½ 4=5 3=5 ¼
3=5 12=25 9=25

0 T 1 0 32=25 24=25 7=25 24=25
H2 ¼ I 2ðuuÞ ¼ þ ¼
0 1 24=25 18=25 24=25 7=25
2 3
1 0 0
H2 ¼ 4 0 7=25 24=25 5
0 24=25 7=25
Since the product H2H1A yields

2 32 3 2 3
1 0 0 14 21 14 14 21 14
H2 H1 A ¼ 4 0 7=25 24=25 54 0 49 14 5 ¼ 4 0 175 70 5
0 24=25 7=25 0 168 77 0 0 35
then H2H1A is the desired upper triangular matrix R and there is no necessity of
further steps. The desired matrix Q results from Q = HT1 HT2 . Therefore
2 32 3
6=7 3=7 2=7 1 0 0
6 76 7
Q ¼ H1 T H2 T ¼ 4 3=7 2=7 6=7 54 0 7=25 24=25 5
2=7 6=7 3=7 0 24=25 7=25
2 3
6=7 69=175 58=175
6 7
¼ 4 3=7 158=175 6=175 5
2=7 6=35 33=35
It is easy to verify that the product QR of the two matrices Q and R determined
above is equal to the given matrix A.
The QR decomposition based on the Givens rotations is shown below.
A Givens rotation is the rotation of a column vector in the plane spanned by two
co-ordinate axes. In order to perform a QR decomposition A = QR, each Givens
rotation is a matrix G which, multiplied on the left by the matrix A, introduces a
zero in the sub-diagonal entries of A, for the purpose of gradually transforming
A into an upper triangular matrix R. The product of the transposes of all these
Givens rotation matrices produces the orthogonal matrix Q.
A Givens rotation matrix is represented by the matrix shown below.
2 3
1 0 0 0
6 .. .. .. .7
6. . . .. 7
6 7
60 cos h sin h 07
6. .. .. .7
G6
6.
. . . .. 7
7
60 sin h 07
6 cos h 7
6 .. .. .. .. 7
4. . . . 5
0 0 0 1
In other words, a Givens rotation matrix G with entries gij results from the
identity matrix I after operating the following substitutions:
gii ¼ cos h gij ¼ sin h

gjj ¼ cos h gji ¼ sin h
This method applies to the case of a QR decomposition as follows.

Let us consider again the matrix
2 3
12 51 4
A4 6 167 68 5
4 24 41
which was decomposed previously. We want to construct a Givens rotation matrix

G1 for the purpose of replacing the a31 entry of A (which entry is at present
a31 = −4) with a zero. By multiplying on the left the first column vector
2 3
12
a1 4 6 5
4
of the matrix A [a1 a2 a3] by the following rotation matrix G1

2 3
1 0 0
G1 4 0 cos h sin h 5
0 sin h cos h
there results the following column vector

2 3
12
4 6 cos h 4 sin h 5
6 sin h 4 cos h
In order for the third entry (−6 sin h − 4 cos h) of this vector to be zero, h must
be

2
h ¼ arctan ¼ 33 :690
3
and consequently cos(−33°.690) = 0.83205 and sin(−33°.690) = −0.55470.

Therefore, G1 is
2 3
1 0 0
G1 4 0 0:83205 0:55470 5
0 0:55470 0:83205
The product G1 A yields a matrix A′, whose entry a′31 is equal to zero, as
follows
2 32 3
1 0 0 12 51 4
A ¼ 40
0
0:83205 0:55470 54 6 167 68 5
20 0:55470 0:83205 4
3 24 41
12 51 4
¼ 4 7:2111 125:64 33:837 5
0 112:60 71:834
Now, starting from A′ indicated above, we find a rotation matrix G2

2 3
cos h sin h 0
G2 4 sin h cos h 05
0 0 1
in order to transform a′21 = 7.2111 into zero. By operating as has been shown
above, we find the following condition
ð sin hÞ 12 þ ðcos hÞ 7:2111 ¼ 0
hence h = 31°.003 and consequently cos h = 0.85714 and sin h = 0.51508. Then
2 3
0:85714 0:51508 0
G2 4 0:51508 0:85714 05
0 0 1
The product A′′ = G2A′ yields a matrix such that its entries a′′21 and a′′31 are
both of them equal to zero:
2 3
14 21 14
A00 4 0 133:96 31:063 5
0 112:60 71:834
Now, starting from A′′ indicated above, let us find a rotation matrix G3
2 3
1 0 0
G3 4 0 cos h sin h 5
0 sin h cos h
which transforms a′′32 = 112.60 into zero. By operating as has been shown above,
we find the following condition
ð sin hÞ 133:96 þ ðcos hÞ 112:60 ¼ 0
hence h = 40°.049, and consequently cos h = 0.76550 and sin h = 0.64344.

Therefore, G3 is
2 3
1 0 0
G3 4 0 0:76550 0:64344 5
0 0:64344 0:76550
The product G3 A′′ yields a matrix R having all of its sub-diagonal entries equal
to zero, as follows
2 3
14 21 14
R4 0 175 70 5
0 0 35
As has been shown above, R results from R = G3G2G1A. Since G1, G2 and G3
are, each of them, orthogonal matrices, then the desired matrix Q of the decom-
position A = QR results from
Q ¼ GT1 GT2 GT3
In the example considered above, the product indicated above yields

2 3
6=7 69=175 58=175
Q 4 3=7 158=175 6=175 5
2=7 6=35 33=35
It is easy to verify that the product QR of the two matrices Q and R determined
above is equal to the given matrix A. Generally speaking, the QR decomposition
A = QR can be used to solve systems of linear algebraic equations. This is because,
by using this decomposition, the system
Ax ¼ b
can be written as follows
QRx ¼ b
Now, since Q is an orthogonal matrix (such that QTQ = I), then multiplying on
the left the two members of the preceding equality by QT yields
R x ¼ QT b
Since R is an upper triangular matrix, then the preceding expression can be

solved for x by back-substitution, as will be shown below. In the specific case of the
least-squares estimation, the QR decomposition A = QR can be used to write the
loss function J described in the preceding paragraph, that is,
J ¼ ðb AxÞT ðb AxÞ
in the way shown below. By multiplying on the left the two members of the
preceding equality by QT and remembering that QTQ = QQT = I, there results
J ¼ ðQT b QT AxÞT ðQT b QT AxÞ ¼ ðd RxÞT ðd RxÞ þ rT r
where the two vectors d and r are partitions of the matrix QTb. The number of their
rows is m for d and n − m for r. The minimum value (rTr) of the loss function J is
reached for Rx = d.
If m is the rank of A, m will also be the rank of R. Consequently, the following

system of linear algebraic equations
2 32 3 2 3
r11 r12 ... r1m x1 d1
6 0 r22 ... r2m 7 6 x 2 7 6 d2 7
6 76 7 6 7
6 .. .. .. 76 .. 7 ¼ 6 .. 7
4 . . ... . 54 . 5 4 . 5
0 0 . . . rmm xm dm
has a unique solution. Since R is an upper triangular matrix, the components x1, x2,
…, xm of x can be computed by back-substitution, that is, beginning from the last
(xm) and going towards the first (x1), as follows
xm ¼ rdmmm
dm1 rðm1Þm xm
xm1 ¼ rðm1Þðm1Þ
..
. !
P
m
xi ¼ r1ii di rij xj ði ¼ m 1; m 2; . . .; 1Þ
j¼i þ 1
In addition to the QR decomposition methods shown above, a singular value

decomposition may be used to solve a least-squares problem. A singular value
decomposition is particularly suited in case of systems of linear algebraic equations
Ax = b having ill-conditioned coefficient matrices A, that is, in case of matrices
A such that small changes in their entries aij result in large changes in the solution
x of the system, as will be shown below. Singular value decomposition is a means
of decomposing a matrix into the product of three simpler matrices, as the sequel
will show.
Any nonzero n m matrix A of rank r > 0 can be decomposed as follows
A ¼ PDQT
that is, decomposed into the product of an n r matrix P with orthonormal col-
umns (such that PTP = I), an r r diagonal matrix D diag(d1, d2, …, dr−1, dr),
whose nonzero entries are supposed to be d1 d2 dr 0, and an
r m matrix QT with orthonormal columns (such that QTQ = I). Let
2 3
d1 0 ... 0
60 d2 ... 0 7
6 7
D 6 .. .. .. 7
4 . . ... . 5
0 0 ... dr
be the diagonal matrix of the singular value decomposition A = PDQT. Its nonzero
entries d1, d2, …, dr, called the singular values of A, are the positive square roots of
the nonzero eigenvalues ki of the Gram matrix K = ATA associated with A, that is,
1
di ¼ ð ki Þ 2 [ 0
(i = 1, 2, …, r). The columns of P, called the left singular vectors, are the nor-
malised eigenvectors of AAT. The columns of Q, called the right singular vectors,
are the normalised eigenvectors of ATA. As has been shown by Abdi [2], singular
vectors come in pairs of one left and one right singular vector corresponding to the
same singular value. They can be computed either separately or as a pair.
When they are computed as a pair (by rewriting equation A = PDQT in another
form, as will be shown below), then it is possible to:
• compute only one (instead of two) decomposition in eigenvectors; and
• prevent a problem which may arise because the normalised eigenvectors of a
matrix are determined up to a factor equal to −1.
Since singular vectors being pairs of vectors must have compatible parities, then
care must be taken of the signs in front of the components of the eigenvectors;
otherwise, the matrices P and Q, if computed separately, might fail to reconstruct
the given matrix A.
Since P and Q are orthogonal matrices, the preceding expression
A ¼ PDQT
can be rewritten as follows
P ¼ AQD1
Consequently, P results from the product of A, Q and D−1. In addition, D−1 is a

diagonal matrix having its nonzero entries equal to the reciprocals of the corre-
sponding entries of D, because D is a diagonal matrix.
It is to be remembered that the eigenvalues of any n n matrix M are the
scalars k1, k2, …, kn which satisfy the following characteristic equation
detðM kIÞ ¼ 0
and that the eigenvectors of M are the corresponding nonzero vectors v such that
ðM kIÞv ¼ 0
If v is an eigenvector of M, then any other nonzero vector w, resulting from the

product between v and a scalar, is also an eigenvector of M.
By the way, the number of the singular values of a matrix M is always equal to
the rank of the matrix itself. Therefore, if an n n matrix M has less than n sin-
gular values, then M is singular.
The condition number k(M) of a non-singular n n matrix M is defined as the
ratio between the largest (d1) and the smallest (dn) of the singular values of M, that
is, k(M) = d1/dn.
A singular matrix M is said to have condition number k(M) equal to infinity; a
matrix M having a very large condition number k(M) is said to be ill-conditioned.
In practice, this condition occurs when the condition number k(M) of the given
matrix M is greater than the reciprocal of the precision of the machine used. In case
of single-precision computations, this occurs typically when k(M) is greater than
107.
An example, taken from Leach [47], of a single value decomposition is given
below. Let

0:96 1:72 0:96 2:28
A AT
2:28 0:96 1:72 0:96
be, respectively, the given matrix to be decomposed and its transpose.

The Gram matrix K = ATA associated with A is

0:96 2:28 0:96 1:72 6:12 3:84
AT A ¼ ¼
1:72 0:96 2:28 0:96 3:84 3:88
The eigenvalues k1 and k2 of K = ATA satisfy the equation
detðK kIÞ ¼ 0
that is,

6:12 k 3:84
det ¼0
3:84 3:88 k
By expanding the determinant, there results
ð6:12 kÞð3:88 kÞ 3:84 3:84 ¼ 0

k2 10k þ 9 ¼ 0
The preceding equation has two roots (sorted in descending order, in the
absolute sense), which are
1
k1 ¼ 5 þ ð25 9Þ2 ¼ 9
1
k2 ¼ 5 ð25 9Þ2 ¼ 1
Therefore, the singular values of A are

1
d1 ¼ ð k1 Þ 2 ¼ 3
1
d2 ¼ ð k2 Þ 2 ¼ 1
Thus, d1 = 3 and d2 = 1 are the nonzero entries of the required matrix D, that is,

3 0
D¼
0 1
D−1, the inverse of D, results immediately from

1 1=3 0
D ¼
0 1
Now, for k1 = 9 and k2 = 1, let us compute the eigenvectors of K = ATA. As

shown above, these eigenvectors become column vectors in a matrix (Q) ordered by
the size of the corresponding eigenvalues. In other words, the eigenvector of the
largest eigenvalue, in the absolute sense, becomes the first (that is, the leftmost)
column, the eigenvector of the next largest eigenvalue becomes the second column,
and so on; the eigenvector of the smallest eigenvalue becomes the last (that is, the
rightmost) column of Q.
(For k = k1 = 9)

6:12 9 3:84 2:88 3:84
AT A kI ¼ ¼
3:84 3:88 9 3:84 5:12
Let a and b be the components of the eigenvector q1 corresponding to k1 = 9.

The components of q1 result from (ATA − kI)q1 = 0. Therefore

2:88 3:84 a 0
¼
3:84 5:12 b 0
2:88a þ 3:84b ¼ 0
3:84a 5:12b ¼ 0
Solving either equation for b yields b = 0.75 a. Therefore

a
q1
0:75 a
Dividing the two components by the Euclidean length ||q1|| = 1.25 a of q1 yields
the normalised eigenvector q1, that is,

0:8
q1
0:6
(For k = k2 = 1)

6:12 1 3:84 5:12 3:84
AT A kI ¼ ¼
3:84 3:88 1 3:84 2:88
The components a and b of the corresponding eigenvector q2 result from

AT A kI q2 ¼ 0
Therefore,

5:12 3:84 a 0
¼
3:84 2:88 b 0
5:12 a þ 3:84 b ¼ 0
3:84 a þ 2:88 b ¼ 0
Solving either equation for b yields b = −4a/3. Therefore

a
q2
43 a
Dividing the two components by the Euclidean length ||q2|| = 5a/3 of q2 yields
the normalised eigenvector q2, that is,

0:6
q2
0:8
Therefore, the required matrix Q is given by Q = [q1q2], that is,

0:8 0:6
Q¼
0:6 0:8
As is easy to verify, the product QTQ is equal to the 2 2 identity matrix I.

Therefore, Q is an orthogonal matrix.
Now, as has been shown above, we compute P = AQD−1. The product QD−1 is

0:8 0:6 1=3 0 0:8=3 0:6
QD1 ¼ ¼
0:6 0:8 0 1 0:2 0:8
The product AQD−1 yields P, as follows

0:96 1:72 0:8=3 0:6 0:6 0:8
P¼ ¼
2:28 0:96 0:2 0:8 0:8 0:6
It is easy to verify that P is an orthogonal matrix and that the product PDQT
yields the given matrix A.
The singular value decomposition can be used for computing the pseudo-inverse
of a matrix, which in turn provides a general method for solving the least-squares
problem in the form Ax = b, as will be shown below. Following Golub and Reinsch
[29], let A be a real n m matrix. An m n matrix X is said to be the
pseudo-inverse of A, if X satisfies the following four properties:
AXA ¼ A
XAX ¼ X
ðAXÞT ¼ AX
ðXAÞT ¼ XA
Let A+ denote the unique solution X. If A = PDQT, then the pseudo-inverse A+

is such that
A þ ¼ Q D þ PT
where
8
<1
ðfor di [ 0Þ
D þ ¼ diagðdiþ Þ diþ ¼ di
:0 ðfor di ¼ 0Þ
Consequently, the pseudo-inverse A+ of a given matrix A can easily be com-

puted as a result of the singular value decomposition of A.
This concept applies to the search for the least-squares solution of an
over-determined system of n algebraic equations in m unknowns, where n > m. Let
Ax ¼ b
be such a system in matrix form. Let
q ¼ b Ax
be the residual vector for some vector x. We want to determine the vector xlsq
corresponding to the least-squares solution of Ax = b, that is, the vector corre-
sponding to the minimum possible residual vector q.
To this end, let us consider the squared Euclidean length
kqk2 ¼ qT q ¼ ðb AxÞT ðb AxÞ ¼ bT b 2 xT AT b þ xT AT Ax
of the residual vector q = Ax − b. To determine the minimum value of ||q||2, we

take the derivative with respect to x and set it to zero. This yields
2 AT b þ 2 AT A x ¼ 0
and hence the following m m system of normal equations
AT A x ¼ AT b
Multiplying both terms of the preceding expression on the left by (ATA)−1 yields
1
x ¼ AT A AT b ¼ A þ b
Therefore, the least-squares solution (x) of Ax = b corresponding to the minimum

value of ||q||2 is given by x = A+ b. Weights can be applied to the original system of
equations Ax = b by using a diagonal matrix W R2 diag(r−2 −2
1 , …, rn ).
In this case, the least-squares solution to be found is
1
x ¼ AT W A AT W b
Bebis [6] gives the following (over-determined) system Ax = b of three equa-

tions for the two unknowns x1 and x2:
11x1 þ 2x2 ¼ 0
2x1 þ 3x2 ¼ 7
2x1 x2 ¼ 5
In matrix form, the preceding equations are as follows

2 3 2 3
11 2
0
4 2 x
3 5 1 ¼ 475
x2
2 1 5
As has been shown above, the least-squares solution (x) results from
x ¼ Aþ b
Therefore, the problem reduces to determining the pseudo-inverse A+ of the

given matrix A. Using the expression written above
1
A þ ¼ AT A AT
we compute ATA as follows

2 3

11 2

11 2 2 4 129 18
AT A ¼ 2 3 5¼
2 3 1 18 14
2 1
Since ATA is a 2 2 matrix, it is easy to compute its inverse (ATA)−1 as

follows. Let M be the following 2 2 matrix

a b
M
c d
Then, the inverse matrix M−1 is given by

d=det ðMÞ b=det ðMÞ
M1 ¼
c=det ðMÞ a=det ðMÞ
where det(M) = ad − bc is the determinant of the given matrix M. In this case,

det AT A ¼ 129 14 18 18 ¼ 1482
Therefore

1 14=1482 18=1482
A A ¼
T
18=1482 129=1482
The pseudo-inverse A+ = (ATA)−1AT results from the product

14=1482 18=1482 11 2 2
18=1482

129=1482 2 3 1
0:079622 0:055331 0:0067476
¼
0:040486 0:28543 0:062753
Consequently, the least-squares solution x = A+b of the given system Ax = b is

2 3

0

0:079622 0:055331 0:0067476 4 5 0:42106
7 ¼
0:040486 0:28543 0:062753 1:6842
5
The same result can also be obtained by computing the singular value decom-
position of the given matrix A, as will be shown below.
The matrix ATA has been computed above. Its eigenvalues result from

det AT A k I ¼ ð129 kÞð14 kÞ 18 18 ¼ 0
The preceding equation, solved for k, yields

1
k1 ¼ 131:75 ðk1 Þ2 ¼ 11:478 1
1 ¼ 0:087121
ðk1 Þ2
1
k2 ¼ 11:248 ðk2 Þ2 ¼ 3:3539 1
1 ¼ 0:29816
ðk2 Þ2
For k = k1 = 131.75, the corresponding normalised eigenvector q1 is

0:98852
q1
0:15111
For k = k2 = 11.248, the corresponding normalised eigenvector q2 is

0:15111
q2
0:98852
The orthogonal matrix Q, having q1 and q2 as its columns, is then

0:98852 0:15111
Q¼
0:15111 0:98852
The matrix D−1 results from

!

1 1 1 0:087121 0
D ¼ diag 1 ; 1 ¼
k21 k22 0 0:29816
The product QD−1 is equal to

0:086121 0:045055
QD1 ¼
0:013165 0:29474
The orthogonal matrix P results from the product AQD−1 and is then
2 3
0:97366 0:093875
P ¼ 4 0:13275 0:97433 5
0:18541 0:20463
Finally, A+ results from
A þ ¼ QD þ PT
where D+ = diag(d+i ) = diag[1/(k1)½, 1/(k2)½] has the same entries as those of D−1
computed above. The product QD+PT yields

0:079623 0:055331 0:0067480
0:040487 0:28543 0:062752
This matrix is to be compared with the matrix which has been computed above
by means of the expression A+ = (ATA)−1AT. The singular value decomposition
applies to the least-squares problem as will be shown below. Let
J ¼ ðb AxÞT ðb AxÞ
be the loss function. Let A = PDQT be the matrix of the least-squares problem
Ax = b. By defining the vectors s and t as follows
s ¼ QT x
t ¼ PT b
the condition (ATA)x = ATb for the minimum value of the loss function can be
written as follows
D2 s ¼ D t
If the matrix of the normal equations is non-singular, then the inverse D−1 of the
D matrix exists and the solution of the least-squares problem Ax = b is
s ¼ D1 t
Remembering the definitions of the vectors s and t and the property of

orthogonality (QT = Q−1) of the matrix Q, the preceding expression becomes
x ¼ Q D1 PT b
or, using the column vectors pi and qi of the matrices, respectively, P and Q,
X
m
1
x¼ pTi b qi
i¼1
di
In order to apply the mathematical methods shown above to a problem of orbit

determination, let us consider again the seven observations of the COBE artificial
satellite performed by Healy [33] on the 6th of November 2000. They are:
Observed time (EST) Right ascension (hh:mm:ss) Declination (°) Offset

17:31:29 21:48:00 −16.3 1
17:32:30 21:41:00 −2.0 0
17:33:30 21:31:45 19.3 0
17:34:30 21:14:00 46.9 3
17:35:29 20:28:00 71.9 3
17:36:30 15:01:00 84.6 2
17:37:30 11:03:00 76.1 0
As the satellite came nearest the cross hairs, time was recorded. The offset in the
last column is an estimate of how close to the cross hairs the satellite came.
We want to construct a function which best fits, in the least-squares sense and in
the time interval given above, the seven observations tabulated above.
Since UTC = EST + 5, then (neglecting the difference between UTC and UT1)
the seven EST times indicated above correspond, respectively, to
UT11 ¼ EST1 þ 5 ¼ 22h 31m 29s

UT12 ¼ EST2 þ 5 ¼ 22h 32m 30s
UT13 ¼ EST3 þ 5 ¼ 22h 33m 30s
UT14 ¼ EST4 þ 5 ¼ 22h 34m 30s
UT15 ¼ EST5 þ 5 ¼ 22h 35m 29s
UT16 ¼ EST6 þ 5 ¼ 22h 36m 30s
UT17 ¼ EST7 þ 5 ¼ 22h 37m 30s
The seven values of right ascension, expressed in degrees, are given below.
ð21 þ 48=60Þ 360=24 ¼ 327:00

ð21 þ 41=60Þ 360=24 ¼ 325:25
ð21 þ 31=60 þ 45=3600Þ 360=24 ¼ 322:9375
ð21 þ 14=60Þ 360=24 ¼ 318:50
ð20 þ 28=60Þ 360=24 ¼ 307:00
ð15 þ 1=60Þ 360=24 ¼ 225:25
ð11 þ 3=60Þ 360=24 ¼ 165:75
In the hypothesis of uncorrelated errors of measurement, the weighting matrix is
W ¼ diagðw11 ; w22 ; . . .; w77 Þ

where the values of the weights w11, w22, …, w77 are preliminarily set, all of them,
to unity. In other words, the starting value of the weighting matrix W is taken equal
to I, where I is the 7 7 identity matrix.
Using a method described by several authors (see, e.g., Refs. [14, 32]), called
iteratively reweighted least squares, we intend to update W iteratively, in such a
way as to give less weight to the more uncertain data points. To this end, we use a
mathematical model to predict the law of variation of the right ascension and
declination with time. Then we use the residuals q1, q2, …, q7 as a measure of
uncertainty. These residuals are the differences between the computed (on the basis
of the mathematical model) and the observed data points. The weights chosen in
each iteration are related to the magnitudes of the residuals computed in the pre-
vious iteration, so that a large residual gives rise to a small weight, as will be shown
below. The preliminary values of the weights along the main diagonal of W are
indicated in the following table.
Obs. No. x Right ascension (°) Declination (°) Preliminary weight

1 −1.0000 327.00 −16.300 1.0000
2 −0.66205 325.25 −2.0000 1.0000
3 −0.32964 322.94 19.300 1.0000
4 0.0027701 318.50 46.900 1.0000
5 0.32964 307.00 71.900 1.0000
6 0.66759 225.25 84.600 1.0000
7 1.0000 165.75 76.100 1.0000
The right ascension (a) and declination (d) of the observed satellite vary with
time t according to functions a(t) and d(t), which are not known a priori.
We approximate these unknown functions, within the time interval t1 t t7,
by means of Chebyshev polynomials T0, T1, …, Tn as follows
aðxÞ ¼ a0 T0 ðxÞ þ a1 T1 ðxÞ þ þ an Tn ðxÞ

dðxÞ ¼ d0 T0 ðxÞ þ d1 T1 ðxÞ þ þ dn Tn ðxÞ
where a0, a1, …, an, d0, d1, …, dn are constant coefficients to be determined, and
T0 ðxÞ ¼ 1
T1 ðxÞ ¼ x
Tn þ 1 ðxÞ ¼ 2x Tn ðxÞ Tn1 ðxÞ ðn 2Þ
are Chebyshev polynomials of the first kind. Supposing that a and d are subject to
errors independently of each other, we want to determine the best (in the weighted
least-squares sense) polynomial approximation for a and d. In other words, we want
to determine the unknown values of the coefficients a0, a1, …, an, d0, d1, …, dn
corresponding to the minimum value of some norm of the residual vector (q). In the
case of the classical least-squares method, this norm is the so-called ‘2 norm, that is,
we search the minimum value of the sum of the squares of the residuals. However,
according to Bube and Langan [14], solutions found by using this method tend to
be very sensitive to data points affected by large errors. By contrast, solutions
coming from searching the minimum value of the ‘p norm (where 1 p < 2) are
less sensitive to errors. The method described in Ref. [14] is a hybrid ‘1/‘2 tech-
nique, by means of which a ‘2 fit (or minimum value of the sum of the squared
residuals) is used for small residuals, and a ‘1 fit (or minimum value of the sum of
the absolute residuals) is used for large residuals.
A smooth transition from the search for the minimum ‘2 norm to the search for
the minimum ‘1 norm is obtained by choosing an appropriate value for a positive
parameter e which results from an estimate of the standard deviation of the resid-
uals, as will be shown below. In the general case of n-degree polynomials, the
model matrix A of the system of linear algebraic equations Ax = b is
2 3
T01 T11 . . . Tn1
6 T02 T12 . . . Tn2 7
6 7
A 6 .. .. .. 7
4 . . ... . 5
T07 T17 . . . Tn7
where the first subscript in the entries Tij of A indicates the degree of the Chebyshev
polynomial, and the second subscript indicates the number of observation. In order
for the system Ax = b to be over-determined with our set of seven observations, the
value of n cannot be greater than five. We set n = 5, which corresponds to a
fifth-degree Chebyshev polynomial approximation to the true unknown functions
a(x) and d(x). In order for the argument x of the Chebyshev polynomials to be
within the interval −1 x 1, we have operated a change of variable (from t to
x) in the second column of the preceding table, as follows
2t ðt7 þ t1 Þ
x¼
t7 t1
where the subscript indicates the number of observation.

As to declination (d), we compute the entries Tij of the 7 6 matrix A by
computing the Chebyshev polynomials up to the fifth degree for the seven obser-
vations indicated above. This yields
2 3
1:0000 1:0000 1:0000 1:0000 1:0000 1:0000
6 1:0000 0:66205 0:12338 0:82542 0:96955 0:45837 7
6 7
6 1:0000 0:32964 0:78268 0:84564 0:22516 0:99409 7
6 7
A6
6 1:0000 0:0027701 0:99998 0:0083101 9:99994 0:013850 7
7
6 1:0000 0:32964 0:78268 0:84564 0:22516 0:99409 7
6 7
4 1:0000 0:66759 0:10865 0:81265 0:97639 0:49101 5
1:0000 1:0000 1:0000 1:0000 1:0000 1:0000
The column vectors x and b contain, respectively, the six unknown coefficients
d0, d1, …, d5 of the fifth-degree approximating polynomial
d ðxÞ ¼ d0 T0 ðxÞ þ d1 T1 ðxÞ þ þ d5 T5 ðxÞ
and the seven observed values of declination. They are

2 3 2 3
d0 16:300
6 d1 7 6 2:0000 7
6 7
x 6 .. 7 b6
4 .. 7
5
4 . 5 .
d5 76:100
As shown above, the weighting matrix W is preliminarily set equal to the 7 7

identity matrix I. Now we form the products ATWA and ATWb. Then, the values
of the coefficients d0, d1, …, d5 (contained in the vector x) of the fifth-degree
approximating polynomial indicated above result from the product of the inverse of
ATWA and ATWb, as follows
1
x ¼ AT W A AT W b
By so doing, according to our preliminary evaluation (with W = I) of the

weights w11, w22, …, w77, the polynomial approximation which best fits, in the
weighted least-squares sense, the observed declinations of the satellite is
d ðxÞ ¼ 39:274 T0 ðxÞ þ 54:564 T1 ðxÞ 8:3776 T2 ðxÞ 9:0385 T3 ðxÞ

0:99597 T4 ðxÞ þ 0:67411 T5 ðxÞ
The approximation shown above makes it possible to compute the values of

declination at the given arguments x1, x2, …, x7. These values are compared with
the observed values of declination; then the residuals q1, q2, …, q7 (computed
minus observed values) are also computed, as will be shown below.
Then, the weighting matrix W is updated by means of some non-negative
weighting function f(qi) of the residuals qi. A possible way to do this is to compute
the new values (which will be placed in the rightmost column of the following
table) of the weights wii, as suggested by Bube and Langan [14]:
1
wii ¼
q 2 14
1þ i

where i = 1, 2, …, 7, and e is a positive constant, called damping parameter, whose

value must be chosen by the solver. By so doing, we search the minimum value of
the following hybrid loss function
X7
q 2 14
JðxÞ ¼ 1þ i 1
i¼1

Bube and Langan suggest to take e roughly equal to 1/1.643 0.6 times the
standard deviation r of the residuals.
In the present case, with reference to the residuals q1, q2, …, q7 given in the
following table (which are computed with W = I), we have
q1 þ q2 þ þ q7
l¼ ¼ 0:30654 106
7
" #12
ðq1 lÞ2 þ ðq2 lÞ2 þ þ ðq7 lÞ2
r¼ ¼ 0:0050870
7
r
¼ ¼ 0:0030961
1:1643
and use this value of e to compute new values of the weights, as shown below.
Obs. No. x Decl. comp. (°) Decl. obs. (°) Residual Weight
1 −1.0000 −16.300 −16.300 0.00042725 0.99530
2 −0.66205 −2.0026 −2.0000 −0.0026114 0.87430
3 −0.32964 19.307 19.300 0.0066071 0.65140
4 0.0027701 46.891 46.900 −0.0089111 0.57289
5 0.32964 71.907 71.900 0.0066452 0.64987
6 0.66759 84.597 84.600 −0.0025940 0.87551
7 1.0000 76.100 76.100 0.00043488 0.99513
When the new values of the weights w11, w22, …, w77 have been computed, the
weighting matrix W is updated as follows W = diag(w11, w22, …, w77), and then
the least-squares problem is solved again with the new weighting matrix, updated as
indicated above. In the present case, the rightmost column of the preceding table,
containing the weights to be used for the next iteration, has been filled as follows
1
w11 ¼ " ¼ 0:99530
#14
0:00042725 2
1þ
0:0030961
1
w22 ¼ " ¼ 0:87430
#14
0:0026114 2
1þ
0:0030961
1
w33 ¼ " ¼ 0:65140
2 #14
0:0066071
1þ
0:0030961
1
w44 ¼ " ¼ 0:57289
#14
0:0089111 2
1þ
0:0030961
1
w55 ¼ " ¼ 0:64987
#14
0:0066452 2
1þ
0:0030961
1
w66 ¼ " ¼ 0:87551
2 #14
0:0025940
1þ
0:0030961
1
w77 ¼ " ¼ 0:99513
#14
0:00043488 2
1þ
0:0030961
and consequently the weighting matrix W is updated as follows
W ¼ diagð0:99530; 0:87430; 0:65140; 0:57289; 0:64987; 0:87551; 0:99513Þ
The iterative process described above comes to an end when the set of weights
computed in a given iteration does not differ, within a specified tolerance, from the
set computed in the preceding iteration. With an accuracy of five significant figures,
two further iterations lead to the following results. Since the weights shown below
are the same, within the accuracy of five significant figures, as those computed in
the previous iteration, we stop here the iterative process.
The final approximating polynomial for declination is then
dðxÞ ¼ 39:274 T0 ðxÞ þ 54:564 T1 ðxÞ 8:3772 T2 ðxÞ 9:0385 T3 ðxÞ

0:99662 T4 ðxÞ þ 0:67411 T5 ðxÞ
1 −1.0000 −16.300 −16.300 0.00042725 0.99530
2 −0.66205 −2.0019 −2.0000 −0.0026114 0.87431
3 −0.32964 19.306 19.300 0.0066071 0.65141
(continued)
(continued)
4 0.0027701 46.890 46.900 −0.0089111 0.57289
5 0.32964 71.906 71.900 0.0066452 0.64987
6 0.66759 84.598 84.600 −0.0025940 0.87552
7 1.0000 76.100 76.100 0.00043488 0.99513
The final set of weights is given in the last column of the table shown above.
Now, we use again the 7 7 identity matrix I as the first estimate of the
weighting matrix W, in order to compute the weighted least-squares best fit for the
observed values of right ascension.
The 7 6 matrix A is the same as that shown above for the case of declination.
The column vectors x and b contain now the six unknown coefficients a0, a1, …, a5
of the approximating polynomial
aðxÞ ¼ a0 T0 ðxÞ þ a1 T1 ðxÞ þ þ a5 T5 ðxÞ
and the seven observed values of right ascension. These vectors are
2 3 2 3
a0 327:00
6 a1 7 6 325:25 7
6 7
x 6 .. 7 b6
4 ... 5
7
4 . 5
a5 165:75
The results found iteratively for right ascension are given below.
First iteration (W = I):
Obs. No. x R.A. computed R.A. observed Residual Weight

(°) (°)
1 −1.0000 326.73 327.00 −0.26715 0.99535
2 −0.66205 326.89 325.25 1.6442 0.87412
3 −0.32964 318.78 322.94 −4.1562 0.65137
4 0.0027701 324.10 318.50 5.6040 0.57293
5 0.32964 302.82 307.00 −4.1806 0.64981
6 0.66759 226.88 225.25 1.6297 0.87572
7 1.0000 165.48 165.75 −0.27402 0.99511
where, again, the rightmost column of the table shown above contains the weights
computed for the next iteration. The value of the damping parameter e results from
the residuals q1, q2, …, q7 (computed minus observed values of right ascension)
given above, as shown below:
q1 þ q2 þ þ q7
l¼ ¼ 0:87193 105
7
" #12
ðq1 lÞ2 þ ðq2 lÞ2 þ þ ðq7 lÞ2
r¼ ¼ 3:1996
7
r
¼ ¼ 1:9474
1:1643
By updating the weighting matrix W (set previously equal to I) by means of the

values contained in the rightmost column of the table shown above, the values of
the coefficients a0, a1, a2, a3, a4, and a5 of the fifth-degree approximating poly-
nomial are determined, as follows
aðxÞ ¼ 278:87 T0 ð xÞ 80:293 T1 ðxÞ 39:000 T2 ðxÞ 10:229 T3 ðxÞ

þ 6:2309 T4 ðxÞ þ 9:8938 T5 ðxÞ
These values make it possible to compute new values of right ascension; these,
in turn, are compared with the observed values of right ascension; then the residuals
q1, q2, …, q7 (computed minus observed values) are computed again.
Third (and last) iteration:
Obs. No. x R.A. comp. (°) R.A. obs. (°) Residual Weight
1 −1.0000 326.83 327.00 −0.26715 0.99535
2 −0.66205 326.43 325.25 1.6442 0.87412
3 −0.32964 318.95 322.94 −4.1562 0.65137
4 0.0027701 324.62 318.50 5.6040 0.57293
5 0.32964 302.97 307.00 −4.1806 0.64981
6 0.66759 226.41 225.25 1.6297 0.87573
7 1.0000 165.58 165.75 −0.27402 0.99511
Since the weights are the same, within the accuracy of five significant figures, as
those computed in the previous iteration, we stop here the iterative procedure.
The final approximating polynomial for right ascension is then
aðxÞ ¼ 278:78 T0 ðxÞ 80:289 T1 ð xÞ 39:209 T2 ðxÞ 10:228 T3 ðxÞ

þ 6:6361 T4 ðxÞ þ 9:8892 T5 ðxÞ
The final set of weights is given in the last column of the table shown above.
2.11 The Kalman Filter 273
2.11 The Kalman Filter
The batch estimation method described in the preceding paragraphs is based on the
least-squares principle. It provides an estimate of the state vector x for a spacecraft
at a given epoch by processing the whole amount of observations in each iteration.
To this end, all the observations which are necessary to determine the spacecraft
orbit must be available before the process of their improvement can take place. This
requirement makes the method of batch estimation less desirable than other
methods in real-time or near-real-time applications, which require a
quasi-continuous update of information to produce an estimate of the state vector
x. As has been shown by Vergez et al. [80], one of these applications is the tracking
of an Earth-orbiting satellite by means of ground stations placed on the surface of
the Earth. Since the number of such stations is limited, an orbiting satellite cannot
be tracked continuously. For example, the following figure, due to the courtesy of
NASA [57], shows the orbit, the US ground stations and their acquisition circles
used to track the Landsat 7 satellite.
A ground station can receive data from a satellite, when the satellite is within the
acquisition circle of the station. When the link between the satellite and a ground
station is lost, then the satellite position at some later time must be predicted, in
order for another ground station to be able to establish a new link.
The tracking data result from measurements (range, azimuth angle, altitude
angle, and their rates of change with time) made on an orbiting spacecraft by
established ground stations. Alternative data are range measurements between two
spacecraft and positions determined by using the Global Positioning System (GPS).
Such data require pre-processing to be put into a form useful for orbit state esti-
mation. To this end, the range and angle data must be converted from a topocentric
Earth-fixed co-ordinate system into a geocentric celestial co-ordinate system.
Another of these applications, according to Conway et al. [19], is the orbit
control of a spacecraft, which requires accurate determination of the spacecraft
position and velocity at a specified time.
One of the methods used for this purpose is the Kalman filter, which is a set of
equations aimed at providing an optimal estimation of either the state vector or,
equivalently, the osculating orbital elements of an orbiting spacecraft from a series
of uncertain observations performed at discrete time-steps.
The Kalman filter is optimal because, in case of a linear system and a Gaussian
distribution of errors, it minimises the mean square error of the estimated quantities.
The Kalman filter is also recursive, because new measurements are processed as
they arrive. The term “filter” comes from the theory of signal processing, where the
information contained in a signal affected by noise must be extracted from the
signal received, by filtering out the noise. The same term is used here, because the
operation of estimating the state vector of a spacecraft from measurements affected
by errors is equivalent to filtering out the errors.
The degree of goodness attained by the Kalman filter in performing this task is
measured by the value of the loss function described in Sect. 2.9. In addition, the
Kalman filter uses the history of the improved measurements to predict the sub-
sequent evolution of the system states. To this end, two types of information are
used by the filter:
• observations coming from appropriate measurement apparatuses (e.g. mea-
surements of range, azimuth angle, altitude angle, and the time rate of change of
the range and the angles, made from established ground stations to an orbiting
spacecraft); and
• a mathematical model of the system under observation, describing how each
state depends on the others and how the measurements depend on the states.
This requires the knowledge of the accuracy of both the measurements and the
mathematical model of the observed system.
When the observed system is an orbiting spacecraft, the analytical model used
for this purpose takes account of the central gravitational force and its perturbations.
Some of the force models used for this purpose have been developed by the North
American Aerospace Defence Command (NORAD), that maintains a catalogue
containing the orbital parameters of about 8000 satellites and space debris com-
puted from radar tracking data. Such models are described in Refs. [36], [54], and
[78]. After computing the forces acting upon the spacecraft by means of the force
model chosen, the state vector can be determined by numerically integrating the
equations of motion with their initial conditions.
The result of the estimation is a predicted orbital state at the time of measure-
ment, as well as the state error covariance matrix and the residuals (computed data
minus observed data). The magnitude of such residuals is expected to decrease with
time as the filter converges to a stable estimate.
The classical Kalman filter has been created for linear systems. However, the
behaviour of an orbiting spacecraft is governed by a nonlinear differential equation.
The necessity of dealing with this and other nonlinear systems has given rise to the
extended Kalman filter, in which the system equations are linearised about the
reference trajectory. Therefore, the Kalman filter used for orbit determination is an
extended Kalman filter, which estimates either the orbital elements or the Cartesian
components of the state vector.
An example, taken from Rojas [62] will illustrate the operations indicated above.
Let x1, x2, …, xn be the results of measurements executed at times, respectively, t1,
t2, …, tn. In case of an orbiting spacecraft, x1, x2, …, xn are just the state vectors (x1,
x2, …, xn) of the spacecraft resulting from n measurements. In the simple case of a
one-dimensional system, x1, x2, …, xn are scalars.
Let us consider for now the simple case. The mean ln of this time series is
x1 þ x2 þ þ xn
ln ¼
n
When a new measurement xn+1 is executed at time tn+1, the mean ln+1 can be
computed again as follows
x1 þ x2 þ þ xn þ xn þ 1
ln þ 1 ¼
nþ1
However, in order to compute the new value (ln+1) of the mean, it is more
convenient to use the old value (ln) of the mean and make a small correction to it
by means of xn+1, as follows
n xn þ 1
ln þ 1 ¼ ln þ ¼ K ðnln þ xn þ 1 Þ
nþ1 n
where K = 1/(n + 1) is called the gain factor. By adding and subtracting ln on the
right-hand side of the equality
ln þ 1 ¼ K ðnln þ xn þ 1 Þ
and remembering the definition K = 1/(n + 1), there results
ln þ 1 ¼ K ðnln þ xn þ 1 Þ
By so doing, the new value (ln+1) of the mean is expressed as a weighted mean
of the old value (ln) of the mean and the new measurement (xn+1). Since we trust
more the old value (ln) of the mean than the new measurement (xn+1), then the
weight of ln is greater than the weight of xn+1.
Not only the mean, but also the variance (also called quadratic standard devi-
ation) of a time series can be computed recursively, as will be shown below.
Let us consider again the results x1, x2, …, xn of n measurements executed at
times, respectively, t1, t2, …, tn. As shown in Sect. 2.9, the variance r2n of this time
series is defined as follows
1X n
r2n ¼ ð x i ln Þ 2
n i¼1
When the result xn+1 of a new measurement comes to the filter, the new variance
r2n+1 is
1 nX þ1 2 1 nX þ1
r2n þ 1 ¼ x i ln þ 1 ¼ ½xi ln K ðxn þ 1 ln Þ2
n þ 1 i¼1 n þ 1 i¼1
where K is the gain factor defined above. By considering separately the last addend
(i = n + 1) and carrying out the sum from 0 to n, the expression
nX
þ1
½xi ln K ðxn þ 1 ln Þ2
i¼1
X
n
ð1 K Þ2 ðxn þ 1 ln Þ2 þ ½xi ln K ðxn þ 1 ln Þ2
i¼1
By expanding the square of the expression within square brackets, there results
X
n
ð1 K Þ2 ðxn þ 1 ln Þ2 þ nK 2 ðxn þ 1 ln Þ2 þ ð x i ln Þ 2
i¼1
X
n
2K ð x i ln Þ ð x n þ 1 l n Þ
i¼1
The last term of the expression written above is zero, because
X
n
ð x i ln Þ ¼ 0
i¼1
Therefore
nX
þ1 h i Xn
½xi ln K ðxn þ 1 ln Þ2 ¼ ðxn þ 1 ln Þ2 ð1 K Þ2 þ nK 2 þ ð x i ln Þ 2
i¼1 i¼1
The definition K = 1/(n + 1) implies that
ð1 K Þ2 þ nK 2 ¼ nK
In addition, the definition of variance
1X n
n i¼1
implies that
X
n
ðxi ln Þ2 ¼ n r2n
i¼1
Therefore
1 h 2 i h i
r2n þ 1 ¼ nrn þ nK ðxn þ 1 ln Þ2 ¼ ð1 K Þ r2n þ K ðxn þ 1 ln Þ2
nþ1
The whole process comprises the following steps to be taken iteratively.

Let x1, x2, …, xn be the results of n measurements.
We compute the mean
x1 þ x2 þ þ xn
ln ¼
n
and the variance
1X n
n i¼1
as has been shown above. Then we compute:

• the gain factor K = 1/(n + 1) every time a new result xn+1 of the measurements
comes to the filter;
• the new value of the mean ln+1 = ln + K(xn+1 − ln);
• a preliminary estimate rn+1*2 of the new standard deviation r2n+1 by means of
rn+1*2 = r2n + K(xn+1 − ln)2; and
• the correct value of the new standard deviation r2n+1 = (1 − K)rn+1*2.
Let us consider now the general case, in which each of the measured values z is
an n-dimensional vector. Let t0 be the initial time of the process. The sequence of
operations performed by the Kalman filter is described below:
(a) at time t0, an initial estimate x−0 of the m-dimensional state vector (including its
uncertainty, expressed by the covariance matrix P−0 associated with x−0 ) is given
to the filter;
(b) a new estimate, relating to a subsequent time t1, is predicted by the filter on the
basis of the mathematical model, and the uncertainty of this estimate is com-
puted as a function of the initial uncertainty and the accuracy of the mathe-
matical model;
(c) observations performed at time t1 with a certain degree of accuracy provide the
filter with new information, which is compared with the information coming
from the predicted estimate and then used to compute a new updated estimate,
relating to t1, and the uncertainty of this new estimate; and
(d) still another estimate, relating to a subsequent time t2, is predicted as in step (b),
but now this estimate is based on the results of step (c).
This cycle, from step (b) to step (c), goes on until the measurements come to an
end. The sequence of prediction (predict a state vector estimate and its covariance
matrix to next time-step) and update (compute the Kalman gain; update the state
vector estimate with measurements; and compute the new covariance matrix of the
updated state vector estimate) is repeated each time new observations arrive.
Since the estimated state vector x−0 (relating to the initial time t0) contains
m scalar random variables, its uncertainty is measured by the covariance matrix P−0
associated with x−0 , that is, by the matrix having, along its main diagonal, the
variances of these scalar random variables; and, in its off-diagonal terms, the
covariances which represent any correlation between pairs of scalar variables. The
simultaneous measurements z1, z2, … (taken at times, respectively, t1, t2, …) are
also n-dimensional vectors processed sequentially in time by the recursive algo-
rithm. At each cycle, the algorithm considers only the new measurement vector and
the values coming from the previous cycle. Therefore, unlike the batch least-squares
algorithm, the Kalman filter algorithm need not store in memory all past
measurements.
The Kalman filter takes an initial estimated state vector and its estimated
covariance matrix, and computes the weights (the Kalman gain) to be used to
combine this estimate with the first state vector coming from the measurement
executed. This yields an updated state vector estimate. Since the Kalman filter
computes an updated state vector estimate by means of the new measurement, then
the covariance matrix must also be changed, in order to take account of the new
information added by the measurement. Therefore, the uncertainty of the state
vector (expressed by the changed covariance matrix) decreases.
Now the Kalman filter must get ready to receive the next state vector coming
from the next measurement. To this end, the Kalman filter must project ahead the
updated state vector estimate and its covariance matrix to the next measurement
time. The state vector is assumed to change with time according to a linear law of
transformation plus a random noise. The predicted state vector estimate can only
follow this linear law of transformation, because the value of the noise is not
known. The uncertainty of this predicted state vector estimate, measured by its
covariance matrix, increases, because the prediction does not take the added noise
into account. This is the last step of the Kalman filter cycle.
Since the state vectors coming from the measurements are recursively processed,
then their uncertainty tends to decrease, because of the increasing amount of the
information carried by each of them. On the other hand, since their uncertainty
increases in the projection step, an equilibrium will be reached, where the uncer-
tainty decrease, which occurs in the update step, is counter-balanced by the
uncertainty increase, which occurs in the projection step.
If there were no random noise in the evolution of the observed system from one
time-step to the next, the uncertainty of the state vector would reduce to zero.
Since this uncertainty changes with time, so does the Kalman gain, which must
therefore be recomputed in each cycle.
We can illustrate now the equations of the extended Kalman filter in the general
case, in which the measured values are n-dimensional vectors. Following
Montenbruck and Gill [53], let z be a one of the batches, that is, one of the
partitions, which make up the whole set of observations. Let x0 be the m-dimen-
sional state vector at the reference epoch t0. Let x−0 be an a priori estimate (forecast)
of x0, as indicated by a superscript minus sign, such that
x
0 ¼ x0 þ Dx0
ref
−
where xref
0 is a common state vector of reference, around which x0 is linearised, and
Dx0 is the increment to be added to x0 to obtain the estimate x−0 .
− ref
Let P−0 be an estimate of the covariance matrix associated with x−0 . The estimates
x0 and P−0 can be obtained as a result of either the processing of previous data
−
batches or initial information.

Let H be the Jacobian H = (∂h/∂x0)ref, that is, the matrix of the partial deriva-
tives of the computed values with respect to the state vector (here the subscript ref
indicates that the partial derivatives are to be evaluated at x0 = xref
0 ) at the reference
epoch t0. The vector-valued function h is the function appearing in the equation
z = h(x0) + e, where e is the vector containing the measurement errors.
The measurement vector z and the a priori estimate x−0 can be used to obtain an
improved estimate Dx+0 (where this improved estimate is indicated by a superscript
plus sign), as follows
Dx0þ ¼ P0þ ððP0 Þ1 þ HT W DzÞ

where Dz = z − h(xref
0 ) is the difference between the actual observations and the
observations predicted on the basis of the reference trajectory, W is the weighting
matrix, and
1
P0 þ ¼ ðP0 Þ1 þ HT H H
is the a posteriori covariance matrix, which takes into account the better knowledge
of x0, resulting from the a priori information combined with the batch z of obser-
vations. The new estimate x+0 is related to the previous estimate x−0 by substituting
the a priori information matrix (P−0 )−1 with the difference
1
ðP0 þ Þ HT W H
between the a posteriori information matrix and the measurement information

matrix. This leads to

Dx0þ ¼ Dx þ T
0 þ P0 H W Dz HDx0
which estimates recursively Dx+0 . In the preceding expression, the matrix
K ¼ P0 þ HT W
(which is called Kalman gain) and the residual vector

ref
q ¼ Dz H Dx
0 ¼ z h x0 H Dx
0
are used to correct the estimates x−0 . The residuals contained in q are computed
taking into account the results z of the measurements, the observations h(xref 0 )
computed according to the reference model, and a linear correction H Dx−0 which
depends on the difference between the estimate x−0 and the reference state xref
0 .
In practice, the a posteriori covariance matrix P+0 is not computed as indicated
above, that is, by using the expression
1
P0 þ ¼ ðP0 Þ1 þ HT W H
because this method requires, at each step, the inversion of an m m matrix, where
m is the dimension of the state vector. Instead, a more convenient way of computing
the a posteriori covariance matrix P+0 is given below
P0 þ ¼ ðI K HÞP0
where the Kalman gain results from

1
K ¼ P0 HT W1 þ H P0 HT
In summary, the recursive estimation algorithm is indicated below. Let the a

priori reference state vector x00 be given together with the covariance matrix P00
associated with it. Let also a series of N measurement batches zI (I = 1, 2, … N) be
given. Then, recursive estimates xI0 of the state vector x0 at epoch and their asso-
ciate covariance matrices PI0 are computed for each batch (I = 1, 2, …, N) by
computing:
• the Kalman gain

I1 T 1
KI ¼ P0 I1 HTI W1
I þ HI P0 HI
• the update of the state vector at epoch

0 I1 0
0 þ KI zI hI x0 HI x0 x0
xI0 ¼ xI1
• the update of its covariance matrix at epoch
P0 I ¼ ðI KI HI ÞxP0 I1
where I designates the m m identity matrix.

The expressions given above refer to an arbitrary number of measurements in
each batch. In practice, each batch comprises only a small number of measurements
taken at a common epoch with possible correlations or even a single scalar
observation.
In case of uncorrelated measurements, it is also possible to process them one at a
time. When this happens, the measurement vector zI is replaced by the corre-
sponding scalar measurement zi; in the same way, the weighting matrix WI is
replaced by the scalar weighting coefficient wi = 1/r2i , the Kalman gain matrix KI is
replaced by the gain vector ki having the same dimension (m) as the state vector x0;
the Jacobian H = (∂h/∂x0)ref is a 1 m matrix (that is, a row vector), and the
products HPHT and Hx are scalar quantities. Consequently, in this case, the three
equations given above become, respectively,
1
ki ¼ P0 i1 HTi r2i þ Hi P0 i1 HTi
0 i1
xi0 ¼ xi1
0 þ ki zi hi x0 Hi x0 x0
0
P0 i ¼ ðI ki Hi ÞP0 i1
In the third of the equations written above, the tensor product ki Hi involves a
column vector (ki) and a row vector (Hi) and yields an m m matrix.
This is because the tensor product of two m-dimensional vectors u and v is by

definition
2 3 2 3
u1 u1 v 1 u2 v 2 ... u1 v m
6 u2 7 6 u2 v 1 u2 v 2 ... u2 v m 7
6 7 6 7
uvT ¼ 6 .. 7½ v1 v2 . . . vm ¼ 6 .. .. .. 7
4 . 5 4 . . ... . 5
um u m v1 um v 2 . . . um v m
The three operations indicated above (computation of the Kalman gain, update
of the state vector at epoch, and update of the covariance matrix of the state vector
at epoch) provide an estimate of the state vector at epoch. In order to complete the
recursion, they are followed by:
• a projection (also called propagation) of the reference state vector xref
0 from
epoch t0 to epoch t1, or (which is the same) a projection of the reference state
ref
vector xi−1 from ti−1 to ti; and
• a projection of the covariance matrix Pi−1 from ti−1 to ti.
−1
To this end, let Ui U(ti, ti−1) = ∂xref
i /∂xi−1 = U(ti, t0)U(ti−1, t0)
ref
denote the
state transition matrix from epoch ti−1 to epoch ti, that is, the matrix whose product
with the state vector xi−1 at time ti−1 gives the state vector xi at time ti, as follows
xi ¼ Ui xi1
The state vector x+i−1 (resulting from data up to and including time ti−1) is used to
predict an a priori state vector x−i at time ti, as follows
þ
x
i ¼ xi þ Ui xi1 xi1
ref ref
and its covariance matrix P−i at time ti, as follows

n T o
P
i ¼ E x
i E xi xi E x i
n þ T T o
þ þ þ
¼ E Ui xi1 E xi1 xi1 E xi1 Ui
þ
¼ Ui Pi1 UTi
+ +
where Pi−1 is the covariance of the state vector xi−1 . At this stage of the recursive
algorithm, the observations zi have not yet been taken into account. Therefore, the
information contained in x−i and in P−i (which are, respectively, the predicted state
vector at time ti and its covariance matrix) is the same as that contained in x+i−1 and
+
in Pi−1 (which are, respectively, the improved state vector at time ti−1 and its
covariance matrix), except for the time to which such quantities refer.
Now, the observations zi at time ti have to be taken into account to update the a
priori information. For this purpose, the residual vector qi is expressed as a function
of quantities related to time ti instead of t0, as follows

qi ¼ zi hi xref
0 H i x
0 x0
ref

¼ zi gi xref
i G i x
i xi
ref
where the function hi(xref

0 ), which is used to predict the observations on the basis of
the reference trajectory, has been replaced by the following function
gi ½ti ; xðti Þ ¼ hi ½ti ; xðt0 Þ
which models the observations as a function of the state vector at the time ti of
measurement. Likewise, the Jacobian Hi has been expressed as follows

@hi @gi @xref
Hi ¼ ref ¼ i
¼ Gi Uðti ; t0 Þ
@x0 @xref
i @xref
0
(where Gi = ∂gi/∂xref
i ) in the equation of the residual vector qi. In the same way, the
Kalman gain is expressed as follows
1
T 1
Ki ¼ P
i Gi Wi þ Gi Pi Gi
T
Levy [49] illustrates the recursive estimation algorithm used in the Kalman filter
by means of a scheme, which is shown in the following figure.
The blocks of this scheme represent the following equations:
Compute Kalman gain ki = Pi−1

0 HTi (r2i + Hi Pi−1
0 HTi )−1
Update state estimates xi0 = xi−1
0 + ki [zi − hi (x00) − Hi (xi−1
0 − x00)]
Compute new covariance Pi0 = (I − ki Hi) Pi−1
0
of updated state estimate
Predict state estimate and x−i = xref
i + Ui (xi−1 − xi−1)
+ ref
−
covariance to next time-step Pi = Ui Pi−1 Ui
+ T
The linearised Kalman filter computes estimates x+i (where i = 1, 2, …, n) of the

state vector xi at measurement times ti and the covariance matrices P+i associated
with these estimates.
Then, the time update phase starts with the integration of the equation of motion
and the variational equations from ti−1 to ti, to obtain the reference state vector xref
i
and the state transition matrix Ui = ∂xref ref
i /∂xi−1. By means of these quantities, the
previous estimate x+i−1 and the associated covariance matrix P+i−1 can be projected
from ti−1 to ti, where ti is the time of the current measurement, as follows
þ
x
i ¼ xi þ Ui xi1 xi
ref ref
P þ
i ¼ Ui Pi1 Ui
T
The measurement update phase starts with the computation of the Kalman gain
Ki, the state vector update x+i , and the covariance matrix update P+i , as follows
1
T 1
Ki ¼ Pi Gi Wi þ Gi Pi Gi
T
ref
xiþ ¼ x
i þ K i z i gi xi G i x
i xi
ref
Piþ ¼ ðI Ki Gi ÞPi
The starting values to be given to the filter are x+0 = xref +

0 and P0 = P0 .
ref
The ordinary Kalman filter described above can be used when the deviations
ref +
between the reference state vector (xi−1 ) and the estimated state vector (xi−1 ) are
small at any time (ti−1). In order to remove this restriction, the extended Kalman
filter has been developed, which derives from the ordinary Kalman filter by
resetting the reference state vector to the estimated state vector at the beginning of
each step. In the successive phase of time update, the current estimated state vector
+
(xi−1 ) is projected ahead from the current epoch (ti−1) to the next (ti), and the
variational equations for the state transition matrix are simultaneously solved. This
yields the predicted state vector x−i at epoch ti
x þ
i ¼ x½ti ; xðti1 Þ ¼ xi1
and the covariance matrix P−i associated with x−i
P þ
i ¼ Ui Pi1 Ui
T
The measurement update phase for the extended Kalman filter differs very little
from that for the ordinary Kalman filter, the only difference being the state vector
update equation, which is more simple in the former case:
1
T 1
Ki ¼ Pi Gi Wi þ Gi Pi Gi
T

xiþ ¼ x
i þ K i z i gi xi
Piþ ¼ ðI Ki Gi ÞP
i
The starting values to be given to the filter are the a priori values of the state
vector (x0 = xapr apr
0 ) and of the associated covariance matrix (P0 = P0 ).
The better performance of the extended Kalman filter is obtained at the cost of a
heavier computational effort in the phase of projection ahead of the state vector and
associated covariance matrix. This is because, when an extended Kalman filter is
used, a separate initial-value problem must be solved by numerical integration for
each measurement to be processed. Montenbruck and Gill [53] point out that
low-order single-step methods (e.g. the fourth-order Runge-Kutta method) are
employed in real-time orbit determination programs based on the extended Kalman
filter.
2.12 Numerical Methods for Kalman Filtering
As is the case with the batch least-squares method, numerical errors may sometimes
impair the performance of the Kalman filter, unless special care is taken in the
computation. Hotop [37] has shown that this is due to the presence of a minus sign
between two matrices in the update expression shown in the preceding paragraph,
that is,
Pkþ ¼ ðI Kk Gk ÞP
k
where the subscript k is the time index. Evaluating this expression on a computer
with single precision (REAL*4) may give rise to negative elements along the main
diagonal of the covariance matrix, which conflicts with the theory of covariance
matrices [37]. As has been shown by Grewal and Kain [30], the covariance matrix
must be symmetric and positive definite (in Sect. 2.9, it has been shown that an
n n real symmetric matrix A is positive definite if the equality zTAz > 0 holds for
all nonzero vectors z with real entries); otherwise that matrix cannot represent valid
statistics for state vector components.
A new formulation of the update expression, with the view of eliminating the
negative diagonal elements, is the so-called Joseph algorithm or stabilised Kalman
algorithm [15]. Without going into the analytical details of the matter, suffice it to
mention that Joseph uses, for the update of the covariance matrix P+k , the following
expression
Pkþ ¼ ðI Kk Gk ÞP T 1 T
k ðI K k G k Þ þ K k Wk K k
which, according to Thornton [69], is mechanised as follows:
W1 ¼ ðI Kk Gk Þ
W2 ¼ W1 P
k
Pkþ ¼ W2 WT1 þ Kk W1 T
k Kk
According to Montenbruck and Gill [53], the Joseph algorithm ensures the
positive definiteness of P+k irrespective of errors in Kk or in (I − Kk Gk).
Another method, due to Bierman [9] and meant to the same effect, computes the
update of the covariance matrix P+k by means of tensor products (see Sect. 2.10), as
follows
V1 ¼ P T
k Gk
P1 ¼ P
k Kk V1
T
V2 ¼ P1 GTk

Pkþ ¼ P1 V2 KTk þ Kk W1 T
k Kk
However, as has been shown by Thornton [69], in either the Joseph or the
Bierman formulation, this method requires more than two times the arithmetic
operations required by the original Kalman method and, in addition, is still sus-
ceptible to numerical errors.
In order to avoid the degradation of the computed covariance matrix P+k , Potter
[61] and Schmidt [63] decompose the covariance matrix as follows
Pk ¼ Sk STk
and derive an algorithm, called square-root method, for recursively computing Sk

instead of Pk. By so doing, symmetric products of triangular factors (Sk and STk ) for
the covariance matrix are updated instead of the covariance matrix itself.
As shown in Sect. 2.9, any m m symmetric positive definite matrix P has a
unique decomposition P = SST, where S is a lower triangular matrix with positive
diagonal entries (sii > 0), and ST is the transpose of S. The decomposition P = SST
can be computed by means of the Cholesky method. The square-root method
guarantees positive definiteness of the computed covariance matrix P+k , because the
2.12 Numerical Methods for Kalman Filtering 287
matrix, if kept in this form, can never have a negative diagonal or become asym-
metric. Therefore, a square-root filter is more stable numerically than a conven-
tional Kalman filter.
Following Born [10], we consider the time update equation
P þ
k ¼ Uk Pk1 Uk
T
concerning the covariance matrix P−k associated with the predicted a priori state
vector x−k at time tk. In order to simplify the notation, the subscripts are dropped, so
that the same equation can be rewritten as follows
P ¼ U P þ UT
Using the decomposition P+ = SST, the preceding equation becomes
P ¼ U S ST UT S ST
where S− = U S. Assuming that the observations are processed one at a time and
that their errors are uncorrelated, the expression written above for the Kalman gain,
that is,
1
T 1
Kk ¼ P
k Gk Wk þ Gk Pk Gk
T
becomes
1
K ¼ S ST GT r2 þ GS ST GT
where r2 is the variance of the observation error. We set by definition

1
a ¼ r2 þ G S ST GT
f ¼ ST GT
−
where a is a scalar and f is an n-dimensional column vector
2 3
f1
6 7
f 4 ... 5
fn
By taking the transposes of the matrices on both sides of the equality
f ¼ ST GT
there results
T T T
f T ¼ ST GT ¼ GT ST ¼ G S
where f −T (the transpose of f −) is the n-dimensional row vector f −T [f −1 , …, f −n ].

This follows from the identities
T
AT ¼A
T
ðA BÞ ¼ BT AT
Therefore, there results

1
a ¼ r2 þ f T f
K ¼ a S f
Consequently, the new covariance matrix P+ (associated with the updated state
estimate x+) can be written as follows
P þ ¼ S þ S þ T ¼ ðI K GÞP ¼ ðI a S f GÞS ST ¼ S ðI a f G S ÞST

¼ S I a f f T ST
Let A− be a matrix defined as follows

A AT ¼ I a f f T
If such a matrix can be found, then the new covariance matrix P+ is expressible
as
P þ ¼ S A AT ST ¼ S þ S þ T
To find A−, we introduce a scalar c such that

A AT ¼ ðI c a f f T ÞðI c a f f T Þ ¼ I a f f T
that is,
I 2 c a f f T þ c2 a2 f f T f f T ¼ I a f f T
Now we define a scalar b such that
b ¼ f T f
It follows that
I 2 c a f f T þ c2 a2 b f f T ¼ I a f f T
that is,

c2 a2 b 2 c a þ a f f T ¼ 0
2
abc 2c þ 1 a f f T ¼ 0
The equality written above is satisfied by the trivial solution a f − f −T = 0. It is

also satisfied by the two values of c for which the expression abc2−2c + 1 vanishes.
These two values are
" #12
1 1 1
c¼

ab ðabÞ2 ab
Remembering the expressions written above
a ¼ ðr2 þ f T f Þ1
b ¼ f T f
there results
1
a¼
r2 þb
hence
1 1 1 ar2
¼ r2 þ b b¼ r2 ¼
a a a
1 a 1 1
¼ ¼
b 1 ar2 ab 1 ar2
Therefore, the two values of c corresponding to abc2 - 2c + 1 = 0 are

1 1
1 1 1 2 1
ðar2 Þ2
c¼
¼
ab a2 b2 ab 1 ar2
If we choose the upper sign (+) in the expression written above, then there
results
1 1
1 þ ðar2 Þ2 1 þ ðar2 Þ2 1
c1 ¼ ¼h ih i¼
1 ar 2 1
1 þ ðar Þ 1 ðar Þ
2 2 2
1
2
1
1 ðar2 Þ2
Otherwise, if we choose the lower sign (−) in the same expression, then there
results
1 1
1 ðar2 Þ2 1 ðar2 Þ2 1
c2 ¼ ¼ h ih i¼
1 ar 2 1 1
1 þ ðar2 Þ2 1 ðar2 Þ2
1
1 þ ðar2 Þ2
In order to prevent the denominator from becoming zero in case of ar2 = 1, we

discard c1 and take
1
c2 ¼ 1
1 þ ðar2 Þ2
as the unique solution of the equation abc2 - 2c + 1 = 0.

Remembering the expressions written above
P þ ¼ S A AT ST ¼ S þ S þ T

A AT ¼ I c a f f T I c a f f T ¼ I a f f T
K ¼ a S f
there results

S þ ¼ S A ¼ S I c a f f T ¼ S c K f T
which is the measurement update for the square-root matrix S.

The resulting computational algorithm is indicated below. In this algorithm, we
use again the subscripts k to indicate the states.
At a given state k = 1, 2, …, n (corresponding to a time tk), the following
quantities are to be specified: x−k (state vector), S−k (square-root matrix associated
with x−k , where S−k is such that P−k = S−k S−T
k ), zk (measurement vector), and Gk
(Jacobian matrix containing the partial derivatives of the measurement vector zk
with respect to the state vector x−k ).
The sequence of computation is indicated below.

(a) f T T
k ¼ Sk Gk
(b) ak ¼ 1= r2 þ f T
k fk
(c) Kk ¼ ak S
k fk
þ

(d) xk ¼ xk þ K k z k G k x

h 1
i k
(e) ck ¼ 1= 1 þ ðak r2 Þ2
(f) Skþ ¼ S T
k ck Kk f k (Sk being such that Pk = Sk Sk )
+ + + +T
(g) Predict, by integrating the differential equation of the reference orbit and
U′ = A U forward in time, the next estimates of the state vector x−k+1 and the
transition matrix U−k+1
(h) Update the square-root matrix S+k and the state vector x+k to k + 1, as follows
(i) S þ
k þ 1 ¼ Uk þ 1 Sk
þ
(j) xk þ 1 ¼ Uk þ 1 xk
(k) Increase k to k + 1 and go to step (a).
The sequence given above is based on the presence of a single observation
(performed at the given time tk) in the measurement vector zk.
Bellantoni and Dodge [7] have extended this method to handle more than a
single observation in zk. At the beginning of the sequence, the initial square-root
matrix S−0 , associated with the initial state vector x−0 , results from the Cholesky
decomposition of the initial covariance matrix P−0 , which in turn must be symmetric
and positive definite.
Throughout the Apollo missions, Potter’s square-root filter has been used for
lunar navigation [48].
On the basis of the square-root method, Carlson [17], Bierman [8], and Thornton
[69] have derived a method which recursively computes an upper triangular
covariance square-root matrix.
This method is based on the upper triangular Cholesky decomposition: any
m m symmetric positive definite matrix P has a unique decomposition
P = UDUT, as will be shown below. For example and without loss of generality, if
P were a 3 3 symmetric and positive definite matrix, then the upper triangular
Cholesky decomposition of P would be
2 3 2 32 32 3
p11 p12 p13 1 u12 u13 d11 0 0 1 0 0
4 p12 p22 p23 5 ¼ 4 0 1 u23 54 0 d22 0 54 u12 1 05
p13 p23 p33 0 0 1 0 0 d33 u13 u23 1
where U is a unit upper triangular matrix and D is a diagonal matrix with

non-negative entries dii (i = 1, 2, …, m).
By so doing, the square-root method shown above is modified as follows
1 1 T
P ¼ UD2 UD2 ¼ SST
where
1
h 1 1 1
i
D2 ¼ diag ðd11 Þ2 ; ðd22 Þ2 ; . . .; ðdmm Þ2
The upper triangular Cholesky decomposition is based on the algorithm shown

below.
Following Bierman [8], for j = m, m − 1, …, 2, we compute the entries djj of the
diagonal matrix D, except the first entry d11, as follows
djj ¼ pjj
Then, for j = m, m − 1, …, 2, we compute the diagonal entries ujj of the upper

triangular matrix U as follows
ujj ¼ 1
Then, for j = m, m − 1, …, 2 and k = 1, 2, ...,j − 1, we compute the off-diagonal

entries ukj of U as follows
pkj
ukj ¼
djj
Then, for k = 1, 2, …, j − 1 and i = 1, 2, …, k, the entries pik of P are destroyed

and replaced by the following new values
pik ¼ pik uij ukj djj
Finally, u11 is set equal to unity, and d11 is set equal to p11.
Bierman [8] has also given a Fortran codification of the algorithm indicated
above. Born [11] has given the following example of UDUT decomposition. Let
2 3
1 2 3
P 42 8 2 5
3 2 14
be the matrix to be decomposed. In case of a 3 3 matrix, we have

2 3 2 32 32 3
p11 p12 p13 1 u12 u13 d11 0 0 1 0 0
4 p12 p22 p23 5 ¼ 4 0 1 u23 54 0 d22 0 54 u12 1 05
p13 p23 p33 0 0 1 0 0 d33 u13 u23 1
which, in the present case, becomes

2 3 2 32 32 3
1 2 3 1 u12 u13 d11 0 0 1 0 0
42 8 2 5 ¼ 40 1 u23 54 0 d22 0 54 u12 1 05
3 2 14 0 0 1 0 0 d33 u13 u23 1
By multiplying the matrices and equating to the given matrix P, there results

u12 ¼ ðp12 p13 p23 =p33 Þ= p22 p223 p33 ¼ ð2 3 2=14Þ= 8 22 14 ¼ 11=54
u13 ¼ p13 =p33 ¼ 3=14
u23 ¼ p23 =p33 ¼ 2=14 ¼ 1=7

d11 ¼ p11 ðp12 p13 p23 =p33 Þ2 = p22 p223 =p33 p213 =p33

¼ 1 ð2 3 2=14Þ2 = 8 22 =14 32 =14 ¼ 7=189
d22 ¼ p22 p223 =p33 ¼ 8 22 =14 ¼ 54=7
d33 ¼ p33 ¼ 14
The matrices U, D, and UT of the decomposition P = UDUT are given below.

2 32 32 3
1 11=54 3=14 7=189 0 0 1 0 0
40 1 1=7 54 0 54=7 0 54 11=54 1 05
0 0 1 0 0 14 3=14 1=7 1
As is easy to verify, the product UDUT reconstructs the given matrix P.

A further example of UDUT decomposition has been given by Kang [46]. Let
2 3
130 186 152 20
6 186 283 230 30 7
P6
4 152
7
230 249 35 5
20 30 35 5
be the matrix to be decomposed. As is easy to verify, P is a square, symmetric, and

positive definite matrix.
By applying the Bierman method shown above, the matrix P is decomposed as
follows
2 3 2 32 32 3
130 186 152 20 1 2 3 4 2 0 0 0 1 0 0 0
6 186 283 230 30 7 60 1 5 67 60 3 0 07 62 1 0 07
6 7¼6 76 76 7
4 152 230 249 35 5 4 0 0 1 7 54 0 0 4 0 54 3 5 1 05
20 30 35 5 0 0 0 1 0 0 0 5 4 6 7 1
Again, the product UDUT on the right-hand side of the preceding equality
reconstructs the given matrix P.
Let us consider now the application of the concepts shown above to the Kalman
filter. Following Kang [46], let
P T
k ¼ Uk Dk Uk
be the UDUT decomposition of the predicted covariance matrix P−k .

Likewise, let
Pkþ ¼ Ukþ Dkþ Ukþ T
be the UDUT decomposition of the updated covariance matrix P+k .

As shown in Sect. 2.11, the updated expression of Pk is
Pkþ ¼ ðI Kk Gk ÞP
k
where I is the identity matrix, Kk is the Kalman gain matrix, and Gk is the Jacobian
matrix containing the partial derivatives of the measurement vector zk with respect
to the state vector x−k . The same expression can also be written
Pkþ ¼ P
k Kk Gk Pk
Now, remembering the following expression (see Sect. 2.11) of the Kalman gain
1
T 1
Kk ¼ P
k Gk Wk þ Gk Pk Gk
T
(where W−1k is the inverse of the weighting matrix Wk) and substituting the
expression of the Kalman gain into P+k = P−k − Kk Gk P−k , there results
1
T 1
Pkþ ¼ P T
k Pk Gk Wk þ Gk Pk Gk Gk P
k
If the expression in parentheses (namely W−1 − T

k + Gk Pk Gk ) were a scalar (s), as
the sequel will show, then the equation written above could be written as follows
Pkþ ¼ P T
k Pk Gk ð1=sÞGk Pk
Now, remembering the UDUT decomposition of the predicted covariance matrix

P−k ,the preceding equation becomes
T T
Pkþ ¼ U T
k Dk Uk ð1=sÞ U T
k Dk Uk Gk Gk Uk Dk Uk
T T T
¼ Uk Dk ð1=sÞ Dk Uk Gk Gk U

k Dk Uk
h T T T T T i T

¼ Uk Dk ð1=sÞ Dk Uk Gk Dk Uk Gk Uk
Since (D−k U−T T − −T T T

k Gk ) (Dk Uk Gk ) is a symmetric and positive definite matrix,
then the term within square brackets is also a symmetric and positive definite
matrix. Therefore, the UDUT decomposition also applies to this term, as follows
h T T T T T i
D
k ð1=sÞ Dk Uk Gk Dk Uk Gk ¼ Uk Dk UT
k
Consequently, the update expression of the covariance matrix Pk becomes

T T
Pkþ ¼ Uk Uk Dk Uk Uk
T T
¼ U
k Uk Dk Uk Uk
T T T
¼ U
k Uk Dk Uk Uk
Setting U+k = U−k U*k and D+k = D*k makes it possible to write
Pkþ ¼ Ukþ Dkþ Ukþ T
which shows that the updated covariance matrix P+k can be decomposed as indicated
above. Now it will be shown that the Kalman gain matrix
1
T 1
Kk ¼ P
k Gk Wk þ Gk Pk Gk
T
can be computed without matrix inversion. To this end, let

2 3
p11 p12 p13 p14
6 p p22 p23 p24 7
P 6 21 7
k 4p p32 p33 p34 5
31
p41 p42 p43 p44 k
be the predicted covariance matrix. Let

1 0 0 0
Gk
0 0 1 0 k
be the Jacobian matrix Gk containing the partial derivatives of the measurement

vector zk with respect to the state vector x−k . Let

r r12
Rk 11
r21 r22 k
be the matrix Rk W−1

k (that is, the inverse of the weighting matrix Wk). Then, the
expression
W1 T
k þ Gk Pk Gk
(which appears in the equation of the Kalman gain) becomes

r11 þ p11 r12 þ p13
r21 þ p31 r22 þ p33 k
As shown in Sect. 2.10, let

a b
M
c d
be a given 2 2 matrix. Then, the inverse matrix M−1 is given by

d=detðMÞ b=detðMÞ
M1 ¼
c=detðMÞ a=detðMÞ
where det(M) = ad − bc is the determinant of the given matrix M.

Consequently, the inverse of the matrix W−1 − T
k + Gk Pk Gk is

ðr22 þ p33 Þ=D ðr12 þ p13 Þ=D
ðr21 þ p31 Þ=D ðr11 þ p11 Þ=D k
where

D ¼ det W1 T
k þ Gk Pk Gk ¼ ðr11 þ p11 Þðr22 þ p33 Þ ðr21 þ p31 Þðr12 þ p13 Þ
When the predicted covariance matrix P−k is a block diagonal and the inverse
weighting matrix Rk W−1 −
k is also diagonal, that is, when the matrices Pk and Rk
are as follows
2 3
p11 p12

6 p21 p22 7 r11 0
Pk ¼ 6

4
7 Rk ¼
p33 p34 5 0 r22 k
p43 p44 k
then the matrix (W−1 − T −1

k + Gk Pk Gk ) is

1=ðr11 þ p11 Þ 0
0 1=ðr22 þ p33 Þ k

1=s1 0
0 1=s2 k
This means that the nonzero entries of the inverted matrix (W−1 − T −1
k + Gk Pk Gk )
result from an operation of scalar division. The Kalman gain matrix
1
T 1
Kk ¼ P
k Gk Wk þ Gk Pk Gk
T
results from the matrix product of P−k GTk by (W−1 − T −1

k + Gk Pk Gk ) , that is,
2 3
p11 p13

6 p21 p23 7
6 7 1=s1 0
4 p31 p33 5 0 1=s2 k
p41 p43 k
This yields the following expression of Kk

2 3
p11 =s1 p13 =s1
6 p21 =s1 p23 =s1 7
Kk ¼ 6
4 p31 =s1
7
p33 =s1 5
p41 =s1 p43 =s1 k
Therefore, the Kalman gain matrix Kk can be computed without matrix inver-
sion. The example given above concerns a 4 4 predicted covariance matrix P−k .
In the general case, the Kalman gain matrix Kk results from
2 3
p11 =s1 p13 =s1 ... p1n =sn1
6 p21 =s1 p23 =s1 ... p2n =sn1 7
6 7
Kk ¼ 6 . .. .. 7
4 .. . ... . 5
pm1 =s1 pm3 =s1 . . . pmn =sn1 k
where
sj ¼ rjj þ Gj P T
k Gk
pij ¼ rjj þ Gij P T
k Gij
It remains to show that the predicted covariance matrix P−k is a block diagonal.
If the updated covariance matrix P+k is a block diagonal, then the predicted
covariance matrix P−k remains block diagonal, because the time update equation
P þ
k ¼ Uk Pk1 Uk
T
(concerning the covariance matrix P−k associated with the predicted a priori state
vector x−k at time tk) can be written as follows
U T þ þ þT
k Dk Uk ¼ Uk ðUk1 Dk1 Uk1 ÞUk
T
Now, since (see above)

þ
T
Pk1 ¼ U T
k1 Uk1 Dk1 Uk1 Uk1
¼ ðU T T
k1 Uk1 ÞDk1 ðUk1 Uk1 Þ
¼ ðU T T T
k1 Uk1 ÞDk1 ðUk1 Uk1 Þ
then
U T T T T
k Dk Uk ¼ Uk ½ðUk1 Uk1 ÞDk1 ðUk1 Uk1 Þ Uk
T
¼ ðUk U
k1 Uk1 ÞDk1 ðUk Uk1 Uk1 Þ
T
− * −
Since (Uk Uk−1 Uk−1 )T is the transpose of (Uk Uk−1 *
Uk−1 ), then the term
ðUk U
k1 Uk1 ÞDk1 ðUk Uk1 Uk1 Þ
T
(that is, the predicted covariance matrix P−k ) is diagonal.

In summary, the preceding analysis has shown that:
• the Kalman gain matrix Kk can be computed without matrix inversion; and
• both the predicted (P−k ) and updated (P+k ) covariance matrices can be decom-
posed as P = UDUT.
The sequence of computation for the square-root method without matrix
inversion is indicated below.
• P þ T
k ¼ Uk Pk1 Uk ¼ Uk Dk Uk
T

t 1
• Kk ¼ P
k Gk Rk þ Gk Pk Gk
T
where
sj ¼ rjj þ Gj P t
k Gk
pij ¼ rjj þ Gij P t

k Gij
2 3
p11 =s1 p13 =s1 ... p1n =sn1
6 p21 =s1 p23 =s1 ... p2n =sn1 7
6 7
Kk ¼ 6 . .. .. 7
4 .. . ... . 5
pm1 =s1 pm3 =s1 . . . pmn =sn1 k
h i
• D
k ð 1=s Þ ðD T T
U
k k G k Þ ðD T T T
U
k k G k Þ ¼ Uk Dk UT
k
• Pkþ ¼ ðU T T T þ þ þT
k Uk ÞDk ðUk Uk Þ ¼ Uk Dk Uk
The updated and predicted state vectors and the transition matrix are computed
as has been shown for the Potter–Schmidt method.
In addition to the numerical methods shown above, we merely mention the
so-called sigma-rho filter, which takes this name from the following expression
p
ij ¼ ri rj qij
given by Grewal and Kain [30] to the entries p−ij of the predicted covariance matrix
P−k . In this expression, r−i is the standard deviation of the ith component, r−j is the
standard deviation of the jth component, and q−ij is the correlation coefficient
between the ith and the jth component of the predicted state vector x−k . It is to be
noted that q is used here for the sole purpose of maintaining the nomenclature used
by Grewal and Kain, and has nothing to do with the residuals. The idea on which
this type of Kalman filter is based is that of updating the standard deviation r and
the correlation coefficients q instead of the predicted covariance matrix. By so
doing, Grewal and Kain express the entries p+ij of the updated covariance matrix P+k
as follows
X
n
riþ rjþ qijþ ¼ r
i rj qij Gs r
s rj qsj
s¼1
The particulars of this method can be found in Ref. [30].
2.13 The Unscented Kalman Filter
As shown in Sect. 2.11, the application of the Kalman filter to nonlinear systems is
sometimes difficult, especially in those of such systems which are highly nonlinear.
The extended Kalman filter, described in Sect. 2.11, linearises all nonlinear models,
in order for the classical Kalman filter to be applied to such nonlinear systems. The
application of extended Kalman filters to nonlinear systems has the following flaws:
1. high instability of the filter, when the assumption of linearity is violated locally;
and
2. difficulty in deriving the Jacobian matrices in most practical cases.
To deal with these systems, there is another class of filters, called sigma-point
Kalman filters, to which class the unscented Kalman filter belongs.
This type of filter, proposed in 1992 by Julier and Uhlmann [42], has since then
been dealt with by several authors (see, e.g., Refs. [70–67]).
According to Julier and Uhlmann, a general nonlinear discrete time system is
represented by the following equations:
xk ¼ f ðxk1 ; uk1 ; vk1 ; k 1Þ

zk ¼ hðxk ; uk ; kÞ þ nk
where xk is the m-dimensional state vector at time-step k, uk−1 is the known input
vector, vk−1 is the q-dimensional state noise process vector due to disturbances and
modelling errors, zk is the observation vector, nk is the measurement noise vector,
and f and h are the nonlinear vector-valued functions (supposed known) repre-
senting the system dynamic model. It is assumed that the noise vectors vk and nk are
zero mean and that, for all i and j, the following equalities hold

E vi vTj ¼ dij Qi

E ni nTj ¼ dij Ri

E vi nTj ¼ 0
where Qi is the covariance matrix of the process noise vk, and Ri is the covariance
matrix of the measurement noise nk. In most practical cases, according to van der
Merwe and Wan [84], the noise terms vk and nk are additive, and the two equations
written above can be simplified as follows
xk ¼ f ðxk1 ; uk1 Þ þ vk1

zk ¼ hðxk Þ þ nk
As shown in Sect. 2.11, the extended Kalman filter consists of:

• a first-order (in terms of Taylor-series expansion) approximation of the non-
linear functions f and h at the current state estimate; and
• the application of the classical Kalman filter to this model approximated to the
first order.
By contrast, the unscented Kalman filter is based on the unscented transform,
which is a numerical procedure for approximating the posterior mean and covari-
ance of a random vector obtained from a nonlinear transformation [67]. The
unscented transform determines a set of sample points, called sigma-points, around
the mean. Such points are chosen so that their mean and covariance should be equal
to, respectively, the mean and covariance of the augmented (see below) state vector,
and then propagated through the nonlinear functions f and h, so as to recover the
mean and covariance of the estimate.
When the noise terms vk and nk are additive, there is no need to augment the state
vector xk and its covariance matrix Pk [68]. By contrast, this need arises when the
observed dynamic system is of the general type indicated above, as will be shown
in detail in Sect. 2.16.
2.13 The Unscented Kalman Filter 301
The initial m 1 state vector x0 has known mean l0 ¼ Eðx0 Þ and covariance
matrix P0 ¼ E½ðx0 l0 Þðx0 l0 ÞT .
For time-step k = 1, 2, …, N − 1 (N being the size of the measurement process),
a set of 2 m + 1 sigma-point vectors vk−1 is computed as follows
v0k1 ¼ xk1 n o i¼0

1
vik1 ¼ xk1 þ ½ðm þ kÞPk1 2 i ¼ 1; 2; . . .; m
n 1
oi
þm
vik1 ¼ xk1 ½ðm þ kÞPk1 2 i ¼ 1; 2; . . .; m
i
where {[(m + k)Pk−1]½}i is the ith column of the matrix Sk−1, which in turn is the
square root of the matrix (m + k)Pk−1. Therefore, there results by definition
ðm þ kÞPk1 ¼ Sk1 STk1
The square-root matrix Sk−1 is to be computed by means of some stable method,

for example, by means of the Cholesky decomposition, as shown in Sect. 2.12. The
set of the sigma-point vectors vk−1
i
(i = 0, 1, …, 2 m) forms the m (2 m + 1)
i
sigma-point matrix Xk−1, whose columns are given below
h i
Xik1 ¼ xk1
1 1
xk1 þ ½ðm þ kÞPk1 2 xk1 ½ðm þ kÞPk1 2
This matrix is such that xk−1 is a column vector, whereas xk−1 + [(m + k)Pk−1]½
and xk−1 − [(m + k)Pk−1]½ are, each of them, a set of m column vectors.
The sigma-point vectors vk−1 are projected ahead from time-step k − 1 to
time-step k, by means of the nonlinear function f, as follows

vik ¼ f vik1 ; uk1 i ¼ 0; 1; . . .; 2m
The projected sigma-point vectors vik (with weights wis for the state vector and wic
for the covariance matrix) are used to compute the predicted state vector x−k and the
predicted covariance matrix P−k of x−k , as follows
X
2m
x
k ¼ wis vik
i¼0
X
2m i T
P
k ¼ wic vik x
k vk x
k þ Qk1
i¼0
where Qk−1 is the covariance matrix of the process noise vector vk−1.
Again, a set of 2m + 1 sigma-point vectors vik (i = 0, 1, …, 2m) is computed

from the predicted state vector x−k and the predicted covariance matrix P−k , as
follows
v0k ¼ x n i¼0
k
12 o
vik ¼ x
k þ ðm þ kÞP
k i ¼ 1; 2; . . .; m
n i 1 o
vikþ m ¼ x
k ðm þ kÞPk 2
i ¼ 1; 2; . . .; m
i
Now, the set of 2m + 1 sigma-point vectors vik is used as the argument of the
nonlinear function h, as follows

cik ¼ h vik i ¼ 0; 1; . . .; 2m
and the sigma-point vectors cik are used (with weights wis for the state vector and wic
for the covariance matrix) to compute the predicted measurement vector z−k and the
predicted covariance matrix P−zz of z−k , as follows
X
2m
z
k ¼ wis cik
i¼0
X
2m i T
P
zz ¼ wic cik z
k ck z
k þ Rk
i¼0
where Rk is the covariance matrix of the measurement noise vector nk.

The cross-correlation matrix P−xz (related to the time-step k) between x−k and z−k is
computed as follows
X
2m i T
P
xz ¼ wic vik x
k ck z
k þ Rk
i¼0
Finally, the Kalman gain Kk, the updated state vector x+k , and the updated
covariance matrix P+k are computed as follows
1
Kk ¼ ðP
xz ÞðPxz Þ
xkþ ¼ x
k þ Kk ðzk zk Þ
Pkþ ¼ P T
k Kk Pzz Kk
According to Julier and Uhlmann [42], the unscented Kalman filter, due to its
better properties of estimation accuracy and ease of implementation in comparison
with the extended Kalman filter, is better suited than the latter in filtering appli-
cations. In quantitative terms, according to Wan and van der Merwe (Refs. [84,
81]), the posterior mean and covariance computed by means of the unscented
2.13 The Unscented Kalman Filter 303
Kalman filter are accurate to the third order, in terms of a Taylor-series expansion,
for all nonlinearities, in case of Gaussian inputs. In case of non-Gaussian inputs, the
approximations are accurate to the second order. In contrast, the extended Kalman
filter provides results approximated to the first order.
2.14 The Square-Root Unscented Kalman Filter
The filter described in the present paragraph has been proposed by van der Merwe
and Wan [84]. It is an improvement of the standard unscented Kalman filter, which
has been described in Sect. 2.13. The difference is due to the fact that the
square-root unscented Kalman filter propagates (that is, projects ahead in time) the
square root Sk of the covariance matrix Pk = Sk STk , instead of the covariance matrix
itself. The reasons for doing this are the same as those discussed in Sect. 2.12,
namely numerical stability and positive definiteness, at each time-step k, of the
covariance matrix Pk of the state vector xk.
Assuming again a system represented by the following equations

we intend to modify the expressions given in Sect. 2.13, so as to take account of the
propagation of the Cholesky factor Sk. The square-root unscented Kalman filter uses
three techniques of linear algebra, namely the QR decomposition, the least-squares
problem, and the Cholesky factor updating. The first two techniques have been
shown in Sect. 2.10. The third will be shown in the last section of the present
paragraph.
Following Terejanu [68], the initial m-dimensional state vector x0 has known
mean l0 ¼ Eðx0 Þ and covariance matrix P0 ¼ E½ðx0 l0 Þðx0 l0 ÞT , whose
Cholesky factor S0 is found as follows
n h io
S0 ¼ chol E ðx0 l0 Þðx0 l0 ÞT
In the preceding expression, the Cholesky factorisation algorithm has been

indicated by means of the function chol(), whose argument is the covariance matrix
P0. The users of MATLAB have just this function at their disposal. Otherwise,
those who have a FORTRAN compiler can use the subroutine schdc of LINPACK
or LAPACK [21], or the subroutine msfd of IBM SSP [39]. Both of the subroutines
named above take a symmetric positive definite matrix A and compute an upper
triangular matrix R (called the Cholesky factor of A) such that A = RTR. For
time-step k within the interval 1, 2, …, kend, the sigma-point vectors vk−1 are
computed as follows
v0k1 ¼ xk1 h i i¼0

1
vik1 ¼ xk1 þ ðm þ kÞ2 Sk1 i ¼ 1; 2; . . .; m
h 1
ii
þm
vik1 ¼ xk1 ðm þ kÞ2 Sk1 i ¼ 1; 2; . . .; m
i
These sigma-point vectors are projected ahead from time-step k − 1 to time-step

k, by means of the nonlinear function f, as follows

vik ¼ f vik1 ; uk1 i ¼ 0; 1; . . .; 2m
predicted Cholesky factor matrix S−k , as follows
P
2m
x
k ¼ wis vik
i¼0nh io
12 i
S
k ¼ qr i ¼ 1; 2; . . .; 2m
1
vk xik ðQk Þ2
wic
h 0 0 0 12 i
S w
k ¼ cholupdate Sk ; vk xk ; sgn wc c
where Qk−1 is the covariance matrix of the process noise vector vk−1. The positive
or negative sign of w0c (that is, w0c > 0 or w0c < 0) determines whether the function
update performs a positive or a negative rank-one update to the Cholesky factori-
sation, as will be shown below. In order to understand the meaning of these
expressions, it is necessary to remember the expression of the predicted covariance
matrix P−k = (S−k )(S−k )T, given in the preceding paragraph, for the standard
unscented Kalman filter, that is,
X
2m i T
P
k ¼ wic vik x
k vk x
k þ Qk1
i¼0
By extracting the first term (corresponding to i = 0) from the sum, this

expression can be written, for i = 1, 2, …, 2 m, as follows
2m
X 12 12 T h i h iT
P x x
1 1
k ¼ wic vik k wic vik k þ Q 2 Q 2
k1 k1
i¼1
0 T
þ wc v0k x
0
k vk x
k
2 12 T 3
h 1 1 i wic vik x 0 T
6 k 7
¼ wic 2 vik x
k Q2 4 1 T 5 þ w0c v0k x
k vk x
k
k1
Q2
k1
2.14 The Square-Root Unscented Kalman Filter 305
1
where Q2k1 is the square-root matrix of the process noise covariance matrix Qk−1.

i 1=2 i 1=
The matrix wc vk x
k Q 2 has m rows and 3m columns. As
k1
1=
=2 1
shown in Sect. 2.10, its 3 m m transpose matrix ½ðwc i Þ 2 ðvik x T
k ÞðQ Þk1 can
be decomposed, by using the QR factorisation, into the product of an orthogonal
3m m matrix Ok and an upper triangular m m matrix (S−k )T, as follows

1 T
1=2 T
wic vik x
k Q =2 ¼ Ok S
k ði ¼ 1; 2; . . .; 2mÞ
k1
The MATLAB function which can be used for this purpose is [Q, R] = qr(A),
where the matrices Q, R, and A are such that A = Q R.
Otherwise, with a FORTRAN compiler, it is possible to use the subroutine sqrdc
of LINPACK or LAPACK [21].
Therefore, the predicted covariance matrix P−k can be expressed as follows
T
P T
k ¼ Sk ðOk Þ ðOk Þ Sk þ w0c ðv0k x T
k Þðvk xk Þ
0
T
¼ S
k Sk þ w0c ðv0k x T
k Þðvk xk Þ
0
In order to include the effect of the term w0c (v0k − x−k )(v0k − x−k )T in the square-root
matrix, it is necessary to perform either a rank-one positive update (if w0c > 0) or a
rank-one negative update (if w0c < 0) to the Cholesky factorisation, as follows
h 0 0 0 12 i
S w
k ¼ cholupdate Sk ; vk xk ; sgn wc c
where sgn(x) is the signum function, defined in Sect. 2.10 (that is, sgn(x) = −1 if
x < 0; sgn(x) = 0 if x = 0; sgn(x) = 1 if x > 0), and cholupdate(R, x, “ + ”) or
cholupdate(R, x, “−”) is the MATLAB function which performs the positive
rank-one update (“+”) or the negative rank-one update (“–”) to the Cholesky
factorisation.
Let R = chol(A) be the Cholesky factor of a given matrix A, as has been shown
above. Then, cholupdate(R, x, “+”) returns the upper triangular Cholesky factor of
A + xxT, where x is a column vector; likewise, cholupdate(R, x, “−”) returns the
upper triangular Cholesky factor of A − xxT.
In the present case, the function cholupdate[S−k , (v0k − x−k ), sgn(w0c )(|w0c |)½] returns
the Cholesky factor of (S−k )(S−k )T + w0c (v0k − x−k )(v0k − x−k )T.
Otherwise, with a FORTRAN compiler, it is possible to use the subroutine schud
(positive update) or schdd (negative update) of LINPACK or LAPACK [21].
By so doing, the predicted covariance matrix can be written P−k = (S−k )(S−k )T.
The same line of reasoning can be followed to compute the predicted covariance
matrix P−zz = (S−z )(S−z )T of z−k and the updated covariance matrix P+k = (S+k )(S+k )T of x+k .
A summary of the equations used in the square-root unscented Kalman filter is

given below. Starting from the initial m-dimensional state vector x0, whose mean
l0 ¼ Eðx0 Þ and covariance matrix P0 ¼ E½ðx0 l0 Þðx0 l0 ÞT are known, the
Cholesky factor S0 of P0 is found as follows
n o
S0 ¼ chol E½ðx0 l0 Þðx0 l0 ÞT
For time-step k within the interval 1, 2, …, kend, the sigma-point vectors vk−1 are
computed as follows
v0k1 ¼ xk1 h i i¼0

1
vik1 ¼ xk1 þ ðm þ kÞ2 Sk1 i ¼ 1; 2; . . .; m
h 1
ii
þm
vik1 ¼ xk1 ðm þ kÞ2 Sk1 i ¼ 1; 2; . . .; m
i
These sigma-point vectors are projected ahead from time-step k − 1 to time-step

k, by means of the nonlinear function f, as follows

vik ¼ f vik1 ; uk1 i ¼ 0; 1; . . .; 2m
predicted Cholesky factor matrix S−k , as follows
X
2m
x
k ¼ wis vik
i¼0
nh 12
io
S
k ¼ qr i ¼ 1; 2; . . .; 2m
1
vik xik
wic ðQk Þ2
h 0 0 0 12 i
S
k ¼ cholupdate S
k ; v k x
k ; sgn wc wc
where Qk−1 is the covariance matrix of the process noise vector vk−1. Again, a set of
2m + 1 sigma-point vectors vik (i = 0, 1, …, 2 m) is computed from the predicted
state vector x−k and the predicted Cholesky factor matrix S−k , as follows
v0k ¼ x
k h i i¼0
S
1
vik ¼ x
kþ ð m þ kÞ 2
k i ¼ 1; 2; . . .; m
h i i
2 S
1
vikþ m ¼ x
k ðm þ k Þ k i ¼ 1; 2; . . .; m
i
Now, the set of 2m + 1 sigma-point vectors vik is used as the argument of the
nonlinear function h, as follows

cik ¼ h vik i ¼ 0; 1; . . .; 2m
and the sigma-point vectors cik are used (with weights wis for the state vector and wic
for the covariance matrix) to compute the predicted measurement vector z−k and the
predicted Cholesky factor matrix S−zz of z−k , as follows
P
2m
z
k ¼ wis cik
i¼0nh io

12 i
S
zz ¼ qr i ¼ 1; 2; . . .; 2m
1
c c
wic i
ðRk Þ 2
h k k 0 0 12 i

Szz ¼ cholupdate Szz ; ck z
0
k ; sgn wc wc
where Rk is the covariance matrix of the measurement noise vector nk.

The cross-correlation matrix P−xz (related to the time-step k) between x−k and z−k is
computed as follows
X
2m i T
P
xz ¼ wic vik x
k ck z
k
i¼0
Finally, the Kalman gain Kk, the updated state vector x+k , and the updated
Cholesky factor matrix S+k are computed as follows
h i
T
Kk ¼ P
xz =ðS zz Þ =S
zz
xkþ ¼ x
k þ Kk ðzk zk Þ
Skþ ¼ cholupdate ðS
k ; Kk Szz ; 1Þ
where / indicates an operation of back-substitution, which is a better alternative to the

operation of matrix inversion used in the standard unscented Kalman filter. Since the
Cholesky factor S−zz is a lower triangular matrix, then the Kalman gain Kk can be
computed by means of two operations of back-substitution in the following expression
Kk ½S T
zz ðSzz Þ ¼ Pxz
Since the quantity U = Kk S−zz, which is the middle argument of the function
S+k = cholupdate(S−k , Kk S−zz, −1), is an m m matrix, then the Cholesky factor S+k
is updated consecutively m times, using the m columns of the matrix U.
As has been shown at the beginning of this paragraph, the advantages of the
square-root unscented Kalman filter over the standard unscented Kalman filter are a
better control of round-off errors and the assurance of positive definiteness of the
covariance matrices Pk associated with the successive state vectors xk. In addition,
as Terejanu [68] points out, the square-root unscented Kalman filter does not
require matrix inversions. As a result of computational experiments performed by
van der Merwe and Wan [84], the square-root unscented Kalman filter is about 20%
faster than its standard counterpart.
Finally, for a better understanding of the matter shown above, the fundamental
concepts on the rank-one positive or negative updates to the Cholesky factorisation
are given below.
Let A {aij} be an n n symmetric positive definite matrix. As shown in

Sect. 2.9, the symmetry of A means that aij = aji, and its positive definiteness means
that vTAv > 0 for all vectors v 6¼ 0. Then, there exists a unique decomposition
A = LLT where L {‘ij} is a lower triangular matrix with positive elements ‘11,
‘22, …, ‘nn along its main diagonal. Let
Ax ¼ b
be a system of linear algebraic equations in matrix form, where A is the matrix

indicated above, and x and b are n-dimensional column vectors.
Following Gill et al. [28], when x has been computed from A and b (see again
Sect. 2.9), it is often necessary to solve a modified system
A x ¼ b
We could of course form the new matrix A* and compute the Cholesky
decomposition A* = L*L*T. However, much less labour is required when L* is
computed directly from L than is required when L* is computed from A*. To this
end, it is necessary to modify the decomposition of A in order to obtain the
decomposition of A*, from which x* can be computed. The modification of
A considered here has the following form
A ¼ A þ azzT
where a is a scalar and z is an n-dimensional vector. Since the quantity a zzT is a

matrix of rank one, then the problem indicated above is called the update of the
Cholesky decomposition of a symmetric positive definite matrix (A) following a
rank-one modification. In particular, for a = ± 1, we consider the positive rank-one
update A* = A + zzT and the negative rank-one update A* = A − zzT.
In case of positive (a > 0) rank-one update, A* is positive definite, and therefore
has a Cholesky decomposition A* = L*L*T.
In case of negative (a < 0) rank-one update, A* may be not positive definite, in
which case it does not possess a Cholesky decomposition. In addition, even when
A* is positive definite, L* may be inaccurately computed, if the modification comes
near to reducing the rank of A [21].
Gill et al. [28] give several methods for modifying matrix factorisations. In
particular, some of these methods can be applied to the Cholesky factorisation. One
of them (called algorithm C1 by the authors cited above) is shown below.
Let us suppose that a given n n symmetric positive definite matrix A (having a
Cholesky factorisation A = LDLT) has been modified by a symmetric matrix of
rank one, as follows
A ¼ A þ a z zT
where a is a scalar and z is an n-dimensional vector.

Starting from the matrices L {‘ij} and D {djj}, which are supposed to be
known, we want to compute two matrices L* {‘*ij} and D* {d*jj}, such that
the modified matrix A* should be equal to
A ¼ L D L T
Supposing that A and A* are, both of them, positive definite (see above), the
recurrence relations for modifying L and D are given below.
1. Define a1 = a, w(1) = z
2. For j = 1, 2, …, n, compute
ðjÞ
pj ¼ w j
djj ¼ djj þ aj p2j
bj ¼ pj aj =djj
aj þ 1 ¼ djj aj =djj
)
ðjÞ
wðrj þ 1Þ ¼ wr pj ‘rj
ðj þ 1Þ r ¼ j þ 1; j þ 2; . . .; n
‘rj ¼ ‘rj þ bj wr
2.15 The Minimax Filter
The Kalman (also called H2) filter is based on the principle of searching the min-
imum variance of the average estimation error for linear systems affected by a
Gaussian noise. However, sometimes the statistical properties of the noise affecting
the system are not known. In such cases, it is necessary to search the minimum of
the worst-case, instead of the average, estimation error.
These limitations have given rise to the minimax (also called H∞) filter. This
type of filter bears this name because it searches the minimum of the maximum
singular value of the transfer function from the noise to the estimation error. Unlike
the Kalman filter, which requires the knowledge of the statistical properties of the
noise affecting the observed process, the minimax filter does not require this
knowledge. Therefore, H∞ filters are more robust (that is, more tolerant or less
sensitive to disturbances and modelling uncertainties) than are Kalman filters.
As shown in Sect. 2.9, the Kalman filter is based on the assumption of a random
distribution, with zero mean value, of the measurement errors. In other words, the
average value of the process noise must be zero, and the average value of the
measurement noise must also be zero. This property must hold not only over the
whole duration of the process, but also at each time instant. Under these conditions,
the Kalman filter leads to the smallest possible quadratic standard deviation (or
2-norm) of the estimation error.
By the way, let e(t) be function representing a scalar signal in the time domain.
Let [a, b] be the interval of definition of the function e(t).
The 2-norm (also written ‘2-norm) of e(t) is defined as follows
0 112
Zb
k e k2 ¼ @ jeðtÞj2 dtA
a
If the integral of the square of the absolute value of e(t) is finite, then the function
e(t) is said to be square integrable.
The infinity-norm (also written ‘∞) of e(t) is defined as follows
kek1 ¼ maxt ½jeðtÞj ða t bÞ
In order to understand why we seek sometimes the minimum of the

infinity-norm instead of the minimum of the 2-norm, let us consider a set of
measured data f(t), which is smooth everywhere over the measurement interval
A a t b, but has outlying values much larger (or much smaller) in a very
narrow sub-interval B contained in A than those outside B. This happens, for
example, when f(t) has either a sharp peak or a sharp notch in B. Let us consider a
polynomial p(t), for example, a fifth-degree polynomial as was the case with the
polynomial considered in Sect. 2.10, defined so that p(t) should be the least-squares
approximation to f(t). In other words, let p(t) be the fifth-degree approximating
polynomial obtained by searching the minimum 2-norm of the standard deviation.
By defining the error as follows: e(t) = f(t) − p(t), the 2-norm of e(t) is
0 112
Zb
k e k2 ¼ @ jeðtÞj2 dtA
a
Since p(t) corresponds to the minimum 2-norm of e(t) = f(t) − p(t), then the
discrepancy e(t) between the measured and the computed data is expected to be in
magnitude near zero outside the sub-interval B and much larger than zero inside B,
just where f(t) has a sharp peak (or notch).
Using the criterion of the minimum 2-norm is tantamount to rejecting outlying
values from a set of measured data, on the grounds that such values appear to be
inconsistent with the other data. This line of conduct is arbitrary and sometimes
unjustified, because it exposes the observer to the risk of ignoring perfectly good
data and the information which they carry with them.
Generally speaking, the question of whether outliers should, or should not, be
retained in a set of measured values is quite difficult. It has been the subject of many
studies since the times of Gauss. The interested reader can find extensive infor-
mation on the matter, for example, in Refs. [3, 24]. For the purpose of the present
paragraph, which describes the minimax filter, it is assumed that this problem has
been solved.
2.15 The Minimax Filter 311
Let us consider once again the problem of estimating the state vector xk of a
system which varies linearly with time. The system in question is governed by the
following equations
xk þ 1 ¼ f ðxk ; uk Þ þ vk ¼ Ak xk þ Bk uk þ vk
zk ¼ hðxk Þ þ nk ¼ Hk xk þ nk
where xk is the m 1 state vector of the system, Ak is an m m matrix, uk is the

r 1 vector of the known input to the system, Bk is an m r matrix, vk is the
random process noise vector, zk is the n 1 measurement vector, Hk is an
n m matrix, and nk is the random measurement noise vector.
We want to estimate the state vector xk+1 on the basis of the measurement vector
zk and our knowledge of the system equations. Following Simon [64], the estimated
state vector xk+1 (that is, x−k+1) is expressed as follows
x
k þ 1 ¼ Ak xk þ Kk ðzk þ 1 Hk Ak xk Þ
where Kk is some gain matrix which is to be determined. If the criterion of the least
2-norm of the estimation error were used, then Kk would be just the Kalman gain,
as has been shown in the preceding paragraphs. By contrast, here we want to
determine Kk on the basis of the least infinity-norm of the estimation error.
Among several possible solutions, Simon cites that one which determines Kk so
that the maximum singular value of the transfer function from the noise to the
estimation error should be less than a given scalar value c. To this end, Simon defines
a loss function J, which is a performance measure of the estimator, as will be shown
below. Let N be the size of the measurement process (in other words, the time-step
k ranges from 0 to N − 1). Let Xk, Vk, and Nk be weighting matrices associated with,
respectively, the estimation error, the process noise, and the measurement noise. Let
X
N 1
P¼ xk x 2
k x
k¼0
X
N 1
Q¼ kvk k2v
k¼0
X
N1
R¼ knk k2n
k¼0
be the weighted squared two-norms of, respectively, the estimation error vector
(xk − x−k ), the process noise vector vk, and the measurement noise vector nk. In the
expressions given above, the notation jjxk jj2x is used to indicate xTk Xk xk, where Xk
is a weighting matrix associated with xk.
Simon [64] defines the loss function J as follows
P
J¼
QþR
As a result of this definition, large values of J correspond to large deviations of

xk from its estimate x−k . The denominator Q + R of the fraction given above may be
considered as the energy of the unknown noise terms; likewise, the numerator P
represents the energy of the estimation error. Let the noise vectors vk and nk be what
they may, a minimax filter is meant to provide a uniformly small estimation error
(xk − x−k ), so that the loss function J should be bounded by a prescribed value [50],
as will be shown below.
In other words, the estimator of a minimax filter tries to determine the estimated
state vector x−k in such a way as to make the value of J as small as possible. To this
end, the filter designer chooses a value of c such that
1
J\
c
The weighting matrices Xk, Vk and Nk must also be chosen by the filter designer
so as to reach desired results.
The state vector estimate x−k which forces J to be less than 1/c is given by Simon
[64] in the following terms:
1
Lk ¼ I Xk Pk þ HTk N1
k Hk Pk
Kk ¼ Ak Pk Lk HTk N1
k

x
k þ 1 ¼ Ak xk þ Bk uk þ Kk zk Hk xk
Pk þ 1 ¼ Ak Pk Lk ATk þ Vk
where Lk is the output matrix, I is the identity matrix, Pk is the error covariance
matrix, and Kk is the minimax gain matrix (whose counterpart, in case of the
Kalman filter, is the Kalman gain matrix).
The condition J < 1/c shows that, in a minimax filter, the ratio of the estimation
error energy to the noise energy is inversely proportional to the value chosen by the
designer for c. In practice, this value cannot be chosen arbitrarily large, because the
mathematical derivation of the minimax filter equation is based on the hypothesis of
an error covariance matrix Pk whose eigenvalues are, in magnitude, less than one. If
we select a value of c which is too large, this condition is not satisfied. In other
words, it is impossible to find an estimator which makes the estimation error
arbitrarily small [64].
The minimax filter equations given above require operations of matrix inversion
at each time-step. In this regard, Simon [64] notes that such operations need not be
performed in practice at each time-step, because the matrices Lk, Pk, and Kk can be
computed off-line. In other words, when the value of c has been properly chosen,
2.15 The Minimax Filter 313
there is no necessity of measurements to recompute the minimax gain matrix Kk at

each step, because the filtering problem can be solved by using a constant value
(K) of the minimax gain matrix. To determine this constant value, Simon [64] notes
that the matrices Kk, computed off-line, approach very quickly a steady-state value,
that is, the matrices Kk converge, after a few time-steps, to a constant matrix K. In a
successive article [65], Simon proposes to compute K by solving the following
simultaneous equations:
K ¼ ðI þ P=cÞ1 PHT
P1 ¼ M1 I=c þ HT H
M ¼ APAT þ I
If the value of c were improperly chosen by the filter designer, then no solution
to these equations would be found.
Another possible method, also suggested by Simon [65], is given below.
1. Form the following 2m 2m matrix

AT AT HT H I=c2
Z¼
AT A þ AT HT H I=c2
2. Find the eigenvectors of Z. Let v1, v2, …, vm be those eigenvectors which

correspond to eigenvalues outside the unit circle.
3. Form the following matrix

X1
½ v1 v 2 . . . vm ¼
X2
where X1 and X2 are m m matrices.

4. Compute M = X2 X−11
Simon [65] notes that this method works only if X1 has an inverse matrix.
Otherwise, the chosen value of c is too large.
As the classical Kalman filter, so the minimax filter has been created for linear
systems. For nonlinear systems (as is the case with the tracking of an orbiting
spacecraft, whose behaviour is governed by a nonlinear differential equation), the
system equations

must be are linearised (by means of a Taylor-series expansion truncated at the first
order) about the reference trajectory. Therefore, the minimax filter is subject to the
same limitations as those which hold with the extended Kalman filter. When one
has to do with nonlinear systems and the robustness (that is, the insensitivity to the
choice of a probability model) of a filter is a concern, it is advisable to use a more

robust unscented Kalman filter than that shown in Sect. 2.13. This type of filter will
be described in the following paragraph.
2.16 A More Robust Unscented Kalman Filter
As shown in Sect. 2.13, a general nonlinear system is governed by the following

equations
xk ¼ f ðxk1 ; uk1 ; vk1 ; k 1Þ

zk ¼ hðxk ; uk ; kÞ þ nk
where xk is the m-dimensional state vector at time-step k, uk−1 is the known input
vector, vk−1 is the process noise vector due to disturbances and modelling errors, zk
is the observation vector, nk is the measurement noise vector, and f and h are the
nonlinear vector-valued functions (supposed known) representing the system
dynamic model.
The initial state vector x0 has known mean vector l0 ¼ Eðx0 Þ and covariance
matrix P0 = E[(x0 − l0)(x0 − l0)T]. The initial state vector x0 is used to construct an
augmented (denoted by the superscript a) ‘-dimensional vector xa0, which results
from concatenating the mean of the true state vector x0 with the mean of the process
noise vector v0 and the mean of the measurement noise vector n0, as shown below:
2 T 3
E x0
xa0 ¼ 4 E vT0 5
E nT0
Likewise, the true covariance matrix P0 of the true state vector x0 is used to
construct the following augmented (denoted by the superscript a) covariance matrix
2 3
P0 0 0
P0 a ¼4 0 Q0 0 5
0 0 R0
where Q0 is the covariance matrix of the process noise vector v0, and R0 is the
covariance matrix of the measurement noise vector n0. Van Zandt [79] has sug-
gested the following improvement to the unscented Kalman filter: adding a ficti-
tious process noise, even if there be none, to the system, for the purpose of taking
account of the uncertainties of the dynamical model used.
2.16 A More Robust Unscented Kalman Filter 315
By so doing, for k = 1, 2, …, N − 1 (N being the size of the measurement

process), the prediction phase and the update phase of the filter are shown below.
Prediction
By using the scaling parameter k (see below), a set of 2‘ + 1 sigma-point vectors
vk−1
i
(i = 0, 1, …, 2‘) is computed from the augmented state vector xk−1 a
and the
a
augmented covariance matrix Pk−1, as follows
v0k1 ¼ xak1 n i¼0

12 o
vik1 ¼ xak1 þ ð‘ þ kÞPak1 i ¼ 1; 2; . . .; ‘
n 12 oi
þ‘
vik1 ¼ xak1 ð‘ þ kÞPak1 i ¼ 1; 2; . . .; ‘
i
a
where {[(‘ + k)Pk−1 ]½}i is the ith column of the matrix Sk−1, which is the square
root of the following matrix
ð‘ þ kÞPak1
so that, by definition, there results
ð‘ þ kÞPak1 ¼ Sk1 STk1
The square-root matrix Sk−1 is to be computed by means of some stable method,

for example, by means of the Cholesky decomposition, as shown in Sect. 2.12.
The set of the sigma-point vectors vk−1
i
(i = 0, 1, …, 2‘) forms the ‘ (2‘ + 1)
i
sigma-point matrix Xk−1, whose columns are given below
h 1 1 i
Xik1 ¼ xak1 xak1 þ ð‘ þ kÞPak1 2 xak1 ð‘ þ kÞPak1 2
a
This matrix is such that xk−1 a
is a column vector, whereas xk−1 + [(‘ + k)Pak −1]½
½
and xk−1 − [(‘ + k)Pk−1] are, each of them, a set of ‘ column vectors.
a a
Now, each of the sigma-point vectors is propagated, that is, projected ahead from
time-step k − 1 to time-step k, by means of the nonlinear function f, as follows

vik ¼ f vik1 ; uk1 ; vik1 ; k 1 ði ¼ 0; 1; . . .; 2‘Þ
where the first argument of the function f is the sigma-point vector due to the true
state vector xk−1, and the third argument is the sigma-point vector due to the process
noise vector vk−1.
Then, these propagated sigma-point vectors vik are weighted (using weights wis
for the state vector and wic for the covariance matrix) and put together in a proper
manner, so as to yield the predicted state vector x−k and the predicted covariance
matrix P−k of x−k , as follows
X
2‘
x
k ¼ wis vik
i¼0
2‘ X
X 2‘ T
P
k ¼ wijc vik vkj
i¼0 j¼0
where the weights wis for the predicted state vector x−k and the weights wic (i = 0, 1,
…, 2‘) for the predicted covariance matrix P−k are as follows
w0s ¼ ‘ þk k ði ¼ 0Þ
w0c ¼ ‘ þk k þ 1 a2 þ b ði ¼ 0Þ
wis ¼ wic ¼ 2ð‘ 1þ kÞ ði 6¼ 0Þ
k ¼ a2 ð‘ þ jÞ ‘
In the expressions given above, the parameters a and j are used to control the
spread of the sigma-points around the mean of the state vector (the values of a and
j are usually set to, respectively, 1 10−3 and 1), and the parameter b is used to
take account of previous knowledge of the distribution of the state vector around its
mean (for a Gaussian distribution, the value of b is set to 2).
Update
In the update phase, each of the sigma-point vectors is used as the argument of
the nonlinear function h, as follows

cik ¼ h vik ; uk ; k þ vik i ¼ 0; 1; . . .; 2‘
where the first argument of the function h is the sigma-point vector due to the
predicted state vector x−k , and the added term is the sigma-point vector due to the
measurement noise vector nk.
The sigma-point vectors cik are weighted (using weights wis for the state vector
and wic for the covariance matrix) and put together in a proper manner, so as to yield
the predicted measurement vector z−k and the predicted covariance matrix P−zz of z−k ,
as follows
X
2‘
z
k ¼ wis cik
i¼0
2‘ X
X 2‘ T
P
zz ¼ wijc cik ckj
i¼0 j¼0
It is to be noted that the matrix P−zz (related to the time-step k) indicated above is
the covariance matrix of the predicted measurement vector z−k , as indicated above.
2.16 A More Robust Unscented Kalman Filter 317
In order to compute the Kalman gain Kk for the unscented Kalman filter, it is
necessary to compute first the cross-correlation matrix P−xz (related to the time-step
k) between x−k and z−k , as follows
2‘ X
X 2‘ T
P
xz ¼ wijc vik ckj
i¼0 j¼0
where vik is the sigma-point vector due to the predicted state vector x−k , as mentioned
above. Hence, the Kalman gain Kk is given by
1
Kk ¼ ðP
xz ÞðPxz Þ
As is the case with the classical Kalman filter, the updated state vector x+k results
from the predicted state vector x−k plus the innovation (zk − z−k ) weighted by the
Kalman gain Kk, as follows
xkþ ¼ x
k þ Kk ðzk zk Þ
and the updated covariance matrix P+k results from the predicted covariance matrix
P−k minus the predicted measurement covariance matrix P−zz (related to k) weighted
by the Kalman gain Kk, as follows
Pkþ ¼ P T
k Kk Pzz Kk
A filter like that described above has been applied by van der Merwe and Wan
[85] to the problem of fusing noisy observations from the Global Positioning
System (GPS), Inertial Measurement Units (IMU), and other available sensors
(such as a barometric altimeter or a magnetic compass). These observations have
been combined with a kinematic or dynamic model of an unmanned aerial vehicle
(specifically, a remotely controlled helicopter). The results obtained by van der
Merwe and Wan by using this type of filter indicate an error reduction of
approximately 30% in both attitude and position estimates relative to the results
coming from an extended Kalman filter [85].
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Chapter 3
The Central Gravitational Force
and Its Perturbations
3.1 The System of Forces Acting on an Earth Satellite
As has been shown in Chap. 2, Sect. 2.1, the two-body scheme applies to the case
of two isolated point-masses which are subject only to their mutual gravitational
attraction. In this case, such masses move, with respect to an inertial reference
system having its origin in their centre of mass, along Keplerian orbits, each of
which is identified by a set of orbital elements.
This ideal scheme is an approximation to the real behaviour of celestial or
artificial bodies, which are subject to other forces, in addition to the central grav-
itational attraction mentioned above. As a result of these additional forces, which
are called perturbations, the actual paths of the bodies subject to them deviate from
the ideal Keplerian orbits. A qualitative description of the main perturbations has
been given in Chap. 2, Sect. 2.1. The purpose of the present chapter is to describe
quantitatively the main perturbing forces, the effects produced by them, and, in
particular, the effects on an artificial satellite revolving around the Earth.
As is well known, the central gravitational attraction is by far the principal force
among those acting on an Earth satellite. This force is conservative, because the
work done by the force of gravity in moving a point-mass along a path between two
points A and B is independent of the path taken for every pair of points (A, B) of
the gravity force field.
In mathematical terms, the gravity force field f, defined over a given region R, is
a conservative vector field, because there is a continuously differentiable scalar
function V, defined over R, such that f = grad(V). The function V is called the
potential of gravitation for the gravity force field f. The central gravitational
attraction between the Earth and an orbiting satellite causes the satellite to move
along a Keplerian orbit around the centre of mass of the Earth. All additional forces
cause modifications to that orbit.
The following table, taken from Blitzer [1], gives the magnitudes (in m/s2) of the
perturbing accelerations acting on Earth satellites orbiting at different altitudes

https://doi.org/10.1007/978-3-319-62220-0_3
324 3 The Central Gravitational Force and Its Perturbations
(in km) in comparison with the central gravity acceleration. Other tables like this
can be found in Montenbruck and Gill [2] and in Walter [3].
It is to be noted that the values of the atmospheric drag and the solar radiation
pressure given in the table refer to a satellite (Vanguard I) whose area-to-mass ratio
is 2.12 10−2 m2/kg and whose mass is 1.47 kg. The area mentioned above is a
reference area for any orbiting satellite (specifically, the cross section of the satellite
perpendicular to the direction of motion).
The meaning of the J coefficients, related to the non-spherical shape of the Earth,
will be explained in the next section. Blitzer [1] points out that the values given for
the luni-solar attraction express not the direct attraction, but rather the effective
perturbing force.
Source 150 km 750 km 1500 km 36000 km

Central gravity 9.35 7.85 6.42 0.22
Earth oblateness J2 30 10−3 20 10−3 14 10−3 160 10−7
J3 0.09 10−3 0.06 10−3 0.04 10−3 0.08 10−7
J4 0.07 10−3 0.04 10−3 0.02 10−3 0.01 10−7
Principal tesseral J22 0.09 10−3 0.07 10−3 0.04 10−3 0.5 10−7
harmonic
Atmospheric drag 3 10−3 10−7 – –
Luni-solar attraction 10−6 10−6 10−6 70 10−7
Solar radiation 10−7 10−7 10−7 10−7
pressure
Ries et al. [4] have considered in particular the perturbing surface forces (such as
atmospheric drag, solar radiation pressure, terrestrial radiation pressure, and thermal
radiation pressure) acting on Earth satellites. Some results (surface forces per unit
mass) found by Ries et al. are given in the following tables.
Satellite h (km) Atmosph. drag (m/s2) Solar rad. press. (m/s2)

−8
ERS-1 770 3 10 7 10−8
Starlette 800–1100 9 10−10 3 10−8
TOPEX/Poseidon 1330 4 10−10 6 10−8
Lageos 5900 1 10−12 4 10−9
Satellite Terrestrial rad. press. (m/s2) Thermal rad. press. (m/s2)

−8
ERS-1 2 10 –
Starlette 8 10−10 –
TOPEX/Poseidon 9 10−9 5 10−9
Lageos 2 10−10 3 10−12
Unless a satellite is perfectly spherical, the forces acting on it generate torques,

which tend to rotate the satellite.
3.2 The Perturbation Due to the Non-spherical Earth 325
3.2 The Perturbation Due to the Non-spherical Earth
The gravity field of the Earth depends on the actual shape and internal distribution
of material of the Earth itself. As to the shape, the Earth has been represented in
Chap. 2, Sect. 2.3, as an ellipsoid of revolution. However, this ellipsoid is an
approximation to the real shape of the Earth. For the purpose of describing this
shape, scientists often refer to the geoid, meaning by this term the equipotential
surface (i.e. the surface characterised by a constant value of the potential of
gravitation V) which would coincide with the mean sea level, if there were no
waves, tides, ocean currents, and other perturbing forces. The National Geodetic
Survey (NGS) of the National Oceanic and Atmospheric Administration (NOAA)
defines the geoid as the equipotential surface of the Earth gravity which best fits, in
a least-squares sense, the global mean sea level [5]. The deviations of the geoid
from the ellipsoid of revolution (see Chap. 2, Sect. 2.3) are elevations (above) or
depressions (below) the ellipsoid. They are shown in the following figure, due to
the courtesy of the NGS [5].
The black solid lines shown above represent equipotential surfaces, each of
which corresponds to a different value of the potential of gravitation V of the Earth.
The blue-dotted line represents the particular equipotential surface corresponding to
the geoid. If the geoid is above the ellipsoid, the value of N is positive. If the geoid
is below the ellipsoid, as shown in the preceding figure, the value of N is negative.
As to the internal distribution of mass, the materials in the Earth interior are of
different composition and not evenly distributed from a place to another. In addi-
tion, this internal distribution is subject to change with time, due to mass transport
and redistribution processes connected with big climate events that cause changes
in the mass of water stored in oceans, continents and atmosphere [6]. In particular,
the constellation of satellites supported by the Satellite Laser Ranging network has
provided accurate measurements reflecting changes in the global mean sea level.
These changes, in turn, result in changes of the Earth gravity force field. Another
space programme meant to the same purpose is the GRACE mission, consisting of
two satellites flying in formation, as fully described in Ref. [7].
According to Cazenave [8], the gravitational potential V at a point P outside the
Earth, due to the heterogeneous mass distribution inside the Earth volume, is
Z
1
V¼G dM
r
M
where G is the gravitational constant, M is the mass of the Earth, and r is the
distance between a mass element dM and the point P.
With reference to the following figure, let us consider first the motion of a
point-mass m, placed in P at time t, subject only to the gravitational attraction of a
perfectly spherical and homogeneous Earth of mass M, radius rE, and centre O.
In practice, an artificial satellite revolving around the Earth can be considered as

a point-mass, because the mass m of the satellite is negligible with respect to the
mass M of the Earth. In addition, the Earth can be considered as a point-mass placed
in O, because it has been assumed perfectly spherical and homogeneous.
Let r OP be the position vector of the point-mass at time t. Let OH be the
projection of OP onto the equatorial plane xy of the Earth. Let AP be the height of
the point-mass with respect to the Earth surface, so that the magnitude r of the
position vector r can be expressed as r = OA + AP.
The position P of the point-mass m at time t can also be expressed in spherical
co-ordinates r, k, and u*, where r is the radius vector, k is the geocentric longitude,
and u* is the geocentric latitude of the point-mass m at time t. In geodesy, it is
customary to use the angle h = 90° − u*, called polar distance or co-latitude,
instead of the geocentric latitude u*. Here, the notation u* (instead of u) is used to
stress the fact that u* is the geocentric (not the geodetic) latitude of P. In accor-
dance with the definition given by Escobal [9], the geocentric latitude u* is the
acute angle (measured in a plane perpendicular to the equatorial plane) which a line
joining the centre of the Earth with a point of the reference ellipsoid makes with the
equatorial plane (u* is in the range −90° u* 90°).
By contrast, the geodetic latitude u is the acute angle (measured in a plane
perpendicular to the equatorial plane) which a line drawn perpendicular to the local
plane tangent to the reference ellipsoid makes with the equatorial plane. In mag-
nitude, the geodetic latitude is always greater than or equal to the geocentric lati-
tude. They have equal values at the poles and at the equator. The maximum
difference between u and u* is approximately 0.19 degrees (that is, 11.5 min of
arc) at a geodetic latitude of 45°6′. The difference between geocentric and geodetic
latitude is shown graphically in the following figure, where the difference between
ellipsoid and sphere is exaggerated for clarity.
The geocentric longitude k is the angle (positive towards east) between the
Greenwich meridian (xz-plane) and the local meridian. This definition holds inde-
pendently of which of the two solids (sphere or ellipsoid) is chosen to represent the
Earth. Therefore, the geographic longitude of a point P of the ellipsoid is also its
geodetic longitude.
The gravitational acceleration (force per unit mass) acting on the point-mass m is
GM r
f ¼
r2 r
where the minus sign is due to the attractive nature of the gravitational force.
Let V be the potential of gravitation, which results from f = grad(V), because f is
a conservative force. Since we use a spherical model as a first approximation to
the actual gravitational field of the Earth, then f acts along the direction of r.
Consequently, V depends only on r, that is, V = V(r). In other terms, the preceding
expression f = −(GM/r2)(r/r) can be written as follows
GM dV
¼
r2 dr
which integrated yields
GM l
V¼ þc ¼ þc
r r
where c is an arbitrary constant, and l = GM is the gravitational parameter of the

Earth. If we choose c = 0 (so that the potential V is zero when the point P is placed
at an infinite distance from O), the preceding expression becomes
l
V¼
r
which is the first-order approximation to the potential V = V(r, k, u*) due to the
actual gravitational force of the Earth. The second-order approximation to the
potential of gravitation V results from considering the Earth as an oblate ellipsoid,
as has been shown in Chap. 2, Sect. 2.3. In this case, as will be shown below, the
potential of gravitation V = V(r, k, u*) is expressible as follows
r 2
l E 1 2

V¼ 1þ C20 3 sin u 1
r r 2
where l is the gravitational parameter of the Earth, r is the radius vector from the
centre of mass of the Earth to the satellite, rE is the equatorial radius of the Earth,
C20 = −0.001082636 is a constant, and u* is the geocentric latitude.
The perturbing term Vp of the total potential of gravitation V results from
subtracting the principal term l/r from the expression given above. This yields

l rE 2 1
Vp ¼ C 20 3 sin2 u 1
r r 2
The perturbing gravitational acceleration a = grad(Vp), due to the non-spherical

shape of the Earth, can be expressed, in a reference system attached to the satellite,
as follows
a ¼ ar ur þ a/ u/ þ ah uh
where ur = r/r is the unit position vector of the satellite, u/ is the unit vector
perpendicular to r, lying in the orbital plane and pointing in the direction of motion
(u/ v > 0), and uh = h/h (where h is the moment-of-momentum vector per unit
mass of the satellite) is the unit vector perpendicular to the orbital plane. In other
words, the radial unit vector ur goes from to the centre of the Earth to the orbiting
body, and is positive outwards. The transverse unit vector u/ also lies in the
instantaneous plane of motion (defined by the instantaneous position and velocity
vectors of the orbiting body) and is positive towards increasing values of true
anomaly /, that is, in the direction of the orbital motion. The bi-normal unit vector
uh is perpendicular to the instantaneous plane of motion and is positive in the
direction of the moment-of-momentum vector h = r v.
Following Merson [10], the three components (respectively, radial, transverse,
and bi-normal) of the perturbing acceleration vector a due to the Earth oblateness
(that is, to C20 only) are expressible as follows
3 l rE 2
ar ¼ C20 ð1 3 sin2 i sin2 uÞ
2 r2 r
3 l rE 2
a/ ¼ 2 C20 sin2 i sinð2uÞ
2r r
3 l rE 2
ah ¼ 2 C20 sinð2iÞ sin u
2r r
where u = x + /, i is the inclination angle of the (instantaneous) orbital plane with

respect to the equatorial plane of the Earth, x is the argument of perigee, and / is
the true anomaly of the satellite at epoch. These expressions of the three compo-
nents (ar, a/, and ah) of the perturbing acceleration vector (a) are introduced into
the Lagrange planetary equations in Gaussian form (see [11]):
2a2 h p i
a0 ¼ ðe sin /Þar þ a/
h r
1
e0 ¼ ½p sin /ar þ ½ðp þ r Þ cos / þ rea/
h
r cosðx þ /Þ
i0 ¼ ah
h

r sinðx þ /Þ
X0 ¼ a/
h sin i

p cos / ðp þ r Þ sin / r sinðx þ /Þ cos i
x0 ¼ ar þ a/ ah
he he h sin i
b
M0 ¼ nþ ½p cos / 2er ar ½ðp þ r Þ sin /a/
ahe
where the prime sign (′) denotes derivatives with respect to time, n = (l/a3)½ is the
mean motion, M = n(t − tperigee) is the mean anomaly, p = a(1 − e2) is the
semi-latus rectum, b = a(1 − e2)½ is the minor semi-axis, and h = (lp)½ is the
magnitude of the moment of momentum per unit mass. By so doing, we see that all
of the orbital elements of a satellite are affected by the perturbing acceleration
vector a due to the non-spherical shape of the Earth.
By the way, when the modified equinoctial orbital elements p, f, g, h, k, and ‘
(see Chap, 1, Sect. 1.10) are used instead of the classical orbital elements in order
to avoid singularities, then the Lagrange planetary equations in Gaussian form
[12, 13] become
"
1 #
0 2p p 2
p ¼ a/
w l

12
p ðw þ 1Þ cos ‘ þ f gðh sin ‘ k cos ‘Þ
f0 ¼ ½sin ‘ar þ a/ ah
l w w

12
p ðw þ 1Þ cos ‘ þ g f ðh sin ‘ k cos ‘Þ
g0 ¼ ½ cos ‘ar þ a/ ah
l w w

12
2
p s cos ‘
h0 ¼ ah
l 2w

12
2
p s sin ‘
k0 ¼ ah
l 2w

2
12

w p h sin ‘ k cos ‘
‘0
1
¼ ðlpÞ 2 þ ah
p l w
where s2 = 1 + h2 + k2, w = 1 + f cos ‘ + g sin ‘, and l is the gravitational

parameter of the Earth.
However, if we consider only the first-order secular (that is, non-periodic) effects
of the oblateness, we see that the ellipsoidal Earth does not affect the major
semi-axis (a), eccentricity (e), or inclination (i) of the orbit of a satellite. In other
words, the variations if the classical elements reduce to (see, e.g. Ref. [9]):
a0 ¼ 0 e 0 ¼ 0 i 0 ¼ 0
3 r 2 cos i
X0 ¼ nC 20
E
2 a ð 1 e2 Þ 2
3 r 2 1 5 cos2 i
x0 ¼ nC 20
E
4 a ð 1 e2 Þ 2
3 r 2 1 3 cos2 i
M 0 ¼ n þ nC 20
E
3
4 a ð1 e2 Þ2
where, again, n = (l/a3)½ is the unperturbed mean motion. Let us consider first the
perturbation induced by the Earth oblateness in the right ascension of the ascending
node (X). In case of polar orbits (i = 90°), X′ = 0, and consequently the direction of
the line of nodes remains constant. For prograde orbits (0 i < 90°), there results
cos i > 0 and then X′ < 0, because C20 has a negative value; consequently, the line
of nodes moves westward (regression of the line of nodes). For retrograde orbits
(90° < i 180°), there results cos i < 0 and C20 < 0; consequently, X′ > 0, which
means that the line of nodes moves eastward.
This fact is used to inject satellites in a helio-synchronous orbit around the Earth.
When, due to some mission requirements, we want the orbital plane of a satellite to
precess about the Earth at the same angular velocity as that at which the Earth
revolves around the Sun, then a helio-synchronous orbit is chosen for that satellite.
In this orbit, the line of nodes rotates eastward by about one degree each day (360°
in a year), in order to follow the revolution of the Earth around the Sun. In
accordance with the fourth of the six equations written above, in order for X′ to be
equal to the mean motion

2p 2p
¼ ¼ 1:991 107 radians/s
T Earth 365:256363 24 3600
of the Earth in its revolution around the Sun, that is, in order for a given orbit to be
helio-synchronous, the following condition must be satisfied
3 l 12 r 2 cos i
1:991 107 ¼
E
C 20
2 a3 a ð 1 e2 Þ 2
The orbit of the Landsat 7 helio-synchronous satellite is shown in the following

figure, which is due to the courtesy of NASA [14].
Using the following constants
l ¼ 398600:4415 km3 =s2

rE ¼ 6378:13630 km
C 20 ¼ 0:001082636
and the following elements
a ¼ 7080:6363 km
e ¼ 0:0001724
relating to the orbit of the Landsat 7 satellite and solving for the cosine (cos i) of the
orbit inclination angle (i) with respect to the equator, there results
2
2 7080:63632 ð1 0:00017242 Þ 1:991 107
7
cos i ¼ ¼ 0:14258
3 398600:4415 ð0:001082636Þ 6378:13632
1
2
hence, i = arccos(−0.14258) = 98°.197.

Let us consider now the perturbation induced by the Earth oblateness in the
argument of perigee (x). Since x′ vanishes for 1 – 5 cos2i = 0, then there are two
values (i = 63°.4 and i = 116°.6) of the orbit inclination for which the line of apsides
does not move. For 0° i < 63°.4 or 116°.6 < i 180°, there results x′ > 0, that
is, the perigee moves in the direction of the motion of the satellite (advance of
perigee). For 63°.4 < i < 116°.6, there results x′ < 0, that is, the perigee moves
opposite to the direction of the motion of the satellite. This fact is used in commu-
nications satellites revolving around the Earth in the so-called Molniya orbit, which
is a highly elliptical (e = 0.722), highly inclined (i = 63°.4) orbit, studied in order for
a satellite to view places of the Earth located at high latitudes.
The Molniya orbit is shown in the preceding figure, which is due to the courtesy
of NASA (15). Since the orbital period is T = 12 h, then a satellite placed in this
orbit passes near the apogee every day and night. As has been shown above, the
value i = 63°.4 of the orbital plane inclination with respect to the equator assures a
constant value for the argument of perigee. This follows from the condition x′ = 0,
which is satisfied if 1 − 5 cos2i = 0, which happens when the inclination angle i is
equal to either 63°.4 or 116°.6.
Now, we come back to the expression
Z
1
V¼G dM
r
M
of the Earth gravitational potential. Let r OP be the position vector of the point P
where the value of the potential V is required, as shown in the following figure.
In this figure, O is the centre of mass of the Earth, s OQ is the position vector of
the elementary mass dM placed in Q inside the Earth, and r OP is the position vector
of the point P outside the Earth where the potential V is to be computed. The elementary
potential dV in P, due to the elementary mass dM in Q, is dV = G(1/|r − s|)dM,
where |r − s| is the magnitude of the vector r − s. Consequently, the potential V in P,
due to the whole mass M of the Earth, is
Z
1
V¼G dM
j r sj
M
With the view of evaluating the integral written above, the expression 1/|r − s| is
expanded in a series of spherical harmonics called Legendre polynomials.
To this end, let c be the angle between r and s, as shown in the preceding figure.
The following relation (cosine rule) holds
s2 s 12
1
jr sj ¼ r 2 þ s2 2r s cos c 2 ¼ r 1 þ 2 cos c
r r
As will be shown at length in Chap. 8, Sect. 8.3, by taking the reciprocal of the
preceding expression and expanding 1/|r − s| in a series of Legendre polynomials,
there results
s2 s 12 1 i
1 1 1X s
¼ 1þ 2 cos c ¼ Pi ðcos cÞ
jr sj r r r r i¼0 r
where Pi(x) is the Legendre polynomial of ith degree. This series converges when
r > s. This is because, after setting for convenience s/r = x, there results
1 1 1 1 12
¼ 1 ¼ 1 ¼ 1 ¼ ð1 zÞ
jr sj ð1 þ x2 2x cos cÞ2 ½1 xð2cos c xÞ2 ð1 zÞ2
where z = x(2 cos c − x) = (s/r)(2 cos c − s/r). Now, we expand f(z) = (1 − z)−½
in a Taylor series around z0 = 0. To this end, we evaluate f(z) and its derivatives up
to the fourth order. This yields
½f ðzÞ0 ¼ ½ð1 zÞ2 0 ¼ 1 ½f 0 ðzÞ0 ¼ ½12 ð1 zÞ2 0 ¼ 12

1 3
½f 00 ðzÞ0 ¼ ½34 ð1 zÞ2 0 ¼ 34 2

5 7
½f ð3Þ ðzÞ0 ¼ ½15
8 ð1 zÞ 0 ¼ 8
15
2 9
½f ð4Þ ðzÞ0 ¼ ½105
16 ð1 zÞ 0 ¼ 16
105
...
Therefore, the expansion of (1 − z)−½ around z0 = 0 is
1 1 3 2 1 15 3 1 105 4
ð1 zÞ2 ¼ 1 þ
1
zþ z þ z þ z þ
2 2! 4 3! 8 4! 16
Hence
1 3
½1 xð2cos c xÞ2 ¼ 1 þ xð2cos c xÞ þ x2 ð2cos c xÞ2
1
2 8
5 3 35 2
þ x ð2cos c xÞ3 þ x ð2cos c xÞ4 þ
16 128
Expanding the powers of the binomial (2 cos c − x) yields
1 3
½1 xð2cos c xÞ2 ¼ 1 þ x cos c x2 þ x2 4 cos2 c 4x cos c þ x2
1
2 8
5 3
þ x 8 cos3 c 12x cos2 c þ 6x2 cos c x3
16
35 4
þ x ð16 cos4 c 32x cos3 c þ 24x2 cos2 c
128
8x3 cos c þ x4 Þ þ
The expression written above, ordered by increasing powers of x, is

1 3 3
½1 xð2cos c xÞ2 ¼ 1 þ x cos c þ x2 þ cos2 c þ x3 cos c
1
2 2 2

5 3 15 2 35 4
þ cos3 c þ x4 cos c þ cos c þ
2 8 4 8
which in turn can be written as follows
½1 xð2cos c xÞ2 ¼ x0 P0 ðcos cÞ þ x1 P1 ðcos cÞ þ x2 P2 ðcos cÞ þ

1
where the nth expression in parentheses, denoted by Pn(cos c), is the Legendre
polynomial of the nth order. If we set for convenience x = cos c, the Legendre
polynomials can be defined recursively as follows
P0 ðxÞ ¼ 1
P1 ðxÞ ¼ x
the polynomials of higher degree P2(x), P3(x), … being given by
ðn þ 1ÞPn þ 1 ðxÞ ¼ ð2n þ 1Þ x Pn ðxÞ n Pn1 ðxÞ
As has been shown above, the first five Legendre polynomials are
1 2
P0 ðxÞ ¼ 1 P1 ðxÞ ¼ x P2 ðxÞ ¼
3x 1
2
1 1
P3 ð xÞ ¼ 5x3 3x P4 ð xÞ ¼ 35x4 30x2 þ 3
2 8
Following Burkard [16], the nth-order Legendre polynomial, Pn(x), is expressible

by means of the Rodrigues formula:
n
1 dn ð x 2 1Þ
Pn ð x Þ ¼
2nn! dxn
Another form of spherical harmonics is given by the associated Legendre

function of the first kind Pnm(x), where m (a non-negative integer smaller than or
equal to n) is the order of the function. According to Montenbruck and Gill [2], the
associated Legendre function of the first kind of degree n and order m is defined as
follows
m dm Pn ð xÞ
Pnm ð xÞ ¼ 1 x2 2
dxm
where Pn(x) is the nth Legendre polynomial defined above. For the sake of com-
pleteness, some authors define the associated first-kind Legendre function as fol-
lows: Pnm(x) = (−1)m [(1 − x2)m/2][dmPn(x)/dxm]. However, we maintain the
definition of Pnm(x) given previously, in accordance with Ref. [2], because of its
consistency with published geopotential coefficients.
The associated Legendre functions of the first kind Pnm(x), also called spherical
harmonics, are solutions of the Legendre associated differential equation

2 00 0 m2
1 x y 2xy þ nðn þ 1Þ y¼0
1 x2
or (in the Sturm–Liouville form)

0 m2
1 x 2 y 0 þ nð n þ 1Þ y¼0
1 x2
where n and m are integers called the degree and order of the function. Here, the
prime sign (′) denotes differentiation with respect to x (that is, y′ dy/dx). This
equation has nonzero solutions that are non-singular over the interval −1 x 1
only when n and m are integers such that 0 m n (or equivalent negative
values). We consider here the important case in which n and m are non-negative
integers. When this condition is satisfied, the general solution of the differential
equation given above can be expressed as follows
yðxÞ ¼ A Pnm ðxÞ þ BQnm ðxÞ
where Pnm(x) and Qnm(x) are called the associated Legendre functions of, respec-
tively, the first kind and the second kind. They are given by
m dm Pn ðxÞ m dm Qn ðxÞ
Pnm ðxÞ ¼ 1 x2 2 Qnm ðxÞ ¼ 1 x2 2
dxm dxm
where Pn(x) is the Legendre polynomial of the nth degree defined above, and

1 1þx
Qn ð xÞ ¼ Pn ð xÞln
2 1x
is the Legendre function of the second kind. The first few of the associated Legendre
functions of the first kind are given in the following table (from Lambeck [17]).
n m Pnm(x) Pnm(sin u*)

0 0 1 1
1 0 x sin u*
(continued)
(continued)
n m Pnm(x) Pnm(sin u*)
1 1 (1 − x2)½ cos u*
2 0 ½(3x2 − 1) ½(3 sin2u* − 1)
2 1 3x(1 − x2)½ 3 cos u* sin u*
2 2 3(1 − x2) 3 cos2u*
3 0 ½(5x3 − 3x) ½(5 sin2u* − 3) sin u*
3 1 ½(1 − x2)½ (15x2 − 3) ½(15 sin2u* − 3) cos u*
3 2 (1 − x2) 15x 15 cos2u* sin u*
3 3 (1 − x2)3/2 15 15 cos3u*
The equation written above

1 i
1 1X s
¼ Pi ðcos cÞ
jr sj r i¼0 r
expresses (apart from G and dM) the potential of gravitation dV due to a point-mass
dM. The expression of dV in terms of Legendre polynomials, integrated over a
continuous distribution of mass, yields the potential of gravitation of the Earth
(having gravitational parameter l = GM and equatorial radius rE) in a point P of
spherical co-ordinates r, k and u*, as follows
l X n n
1 X
r
Vðr; k; u Þ ¼ ½C nm cosðmkÞ þ Snm sinðmkÞPnm ðsin u Þ
E
r n¼0 m¼0
r
where Cnm and Snm are the Stokes coefficients, whose values, determined by ter-
restrial means and observations of satellite orbits [18], depend on the mass distri-
bution within the Earth; and Pnm(x) is the associated first-kind Legendre function of
degree n and order m, according to the definition given above. In terms of spherical
co-ordinates r, k, and h, the preceding expression becomes
l X
1 X n n
rE
Vðr; k; hÞ ¼ ½Cnm cosðmkÞ þ Snm sinðmkÞPnm ðcos hÞ
r n¼0 m¼0 r
in accordance with the expressions given in Refs. [19, 20]. The following figure
(courtesy of NOAA [16]) illustrates the first-kind Legendre functions.
The Stokes coefficients with m = 0 are called zonal coefficients, because they
determine the part of the geopotential which depends only on latitude. Among
them, the Sn0 coefficients are all equal to zero, and the remaining Cn0 coefficients
are also indicated with −Jn, so that Cn0 = −Jn. In addition, there results C00 = 1 and
C10 = C11 = S11 = 0, because the origin O of the reference system is the centre of
mass of the Earth. Therefore, the potential of gravitation V has no terms of degree
n = 1 on the right-hand side of the preceding expressions.
The other Stokes coefficients are called tesseral coefficients (for m < n) and
sectorial coefficients (for m = n). This explains the meaning of the J2, J3, J4, and J22
coefficients shown in the first table (from Ref. [1]) of Sect. 3.1. Since Blitzer calls
J22 the principal tesseral (instead of sectorial) harmonic coefficient, the notation and
name given by Blitzer have been kept in the table.
The following table (courtesy of the University of Texas at Austin, Centre of
Space Research [21]) gives a subset of the Stokes coefficients Cnm and Snm for the
JGM 3 gravity model, which in turn is described in Ref. [22].
n m Cnm Snm
2 0 −0.10826360229840D-02 0.0
3 0 0.25324353457544D-05 0.0
4 0 0.16193312050719D-05 0.0
5 0 0.22771610163688D-06 0.0
6 0 −0.53964849049834D-06 0.0
7 0 0.35136844210318D-06 0.0
8 0 0.20251871520885D-06 0.0
2 1 −0.24140000522221D-09 0.15430999737844D-08
3 1 0.21927988018965D-05 0.26801189379726D-06
4 1 −0.50872530365024D-06 −0.44945993508117D-06
5 1 −0.53716510187662D-07 −0.80663463828530D-07
6 1 −0.59877976856303D-07 0.21164664354382D-07
7 1 0.20514872797672D-06 0.69369893525908D-07
(continued)
(continued)
n m Cnm Snm
8 1 0.16034587141379D-07 0.40199781599510D-07
2 2 0.15745360427672D-05 −0.90386807301869D-06
3 2 0.30901604455583D-06 −0.21140239785975D-06
4 2 0.78412230752366D-07 0.14815545694714D-06
5 2 0.10559053538674D-06 −0.52326723987632D-07
6 2 0.60120988437373D-08 −0.46503948132217D-07
7 2 0.32844904836492D-07 0.92823143885084D-08
8 2 0.65765423316743D-08 0.53813164055056D-08
3 3 0.10055885741455D-06 0.19720132389889D-06
4 3 0.59215743214072D-07 −0.12011291831397D-07
5 3 −0.14926153867389D-07 −0.71008771406986D-08
6 3 0.11822664115915D-08 0.18431336880625D-09
7 3 0.35285405191512D-08 −0.30611502382788D-08
8 3 −0.19463581555399D-09 −0.87235195047605D-09
4 4 −0.39823957404129D-08 0.65256058113396D-08
5 4 −0.22979123502681D-08 0.38730050770804D-09
6 4 −0.32641389117891D-09 −0.17844913348882D-08
7 4 −0.58511949148624D-09 −0.26361822157867D-09
8 4 −0.31893580211856D-09 0.91177355887255D-10
5 5 0.43047675045029D-09 −0.16482039468636D-08
6 5 −0.21557711513900D-09 −0.43291816989540D-09
7 5 0.58184856030873D-12 0.63972526639235D-11
8 5 −0.46151734306628D-11 0.16125208346784D-10
6 6 0.22136925556741D-11 −0.55277122205966D-10
7 6 −0.24907176820596D-10 0.10534878629266D-10
8 6 −0.18393642697634D-11 0.86277431674150D-11
7 7 0.25590780149873D-13 0.44759834144751D-12
8 7 0.34297618184624D-12 0.38147656686685D-12
8 8 −0.15803322891725D-12 0.15353381397148D-12
The perturbing term Vp of the total potential of gravitation V of the Earth (so that
V = l/r + Vp) is
l X
1 X n n
rE
Vp ðr; k; hÞ ¼ ½Cnm cosðmkÞ þ Snm sinðmkÞPnm ðcos hÞ
r n¼2 m¼0 r
It is to be observed that Vp also depends on the gravitational parameter l and

equatorial radius rE of the Earth. Therefore, care must be taken in choosing con-
sistent values of l, rE, Cnm, and Snm. With reference to the JGM 3 gravity model
cited above, the values of the gravitational parameter and equatorial radius of the
Earth are, respectively, l = 398600.4415 km3/s2 and rE = 6378.13630 km.
Since the values of the Stokes coefficients Cnm and Snm differ from one another
by several orders of magnitude, it is convenient to introduce a normalisation of such
coefficients which makes them readily comparable in numerical operations.
By applying a scale factor Nnm which depends on the degree n and on the order
m of the associated Legendre function, the normalised Stokes coefficients Cnm and
Snm (over-lined) are used in place of C nm and Snm . The following table (from Ref.
[22]) gives the same subset of the normalised Stokes coefficients C nm and Snm for
the JGM 3 gravity model as that given in the previous table.
n m C nm Snm
2 0 −484.165368D-06 0.0
3 0 0.957171D-06 0.0
4 0 0.539777D-06 0.0
5 0 0.068659D-06 0.0
6 0 −0.149672D-06 0.0
7 0 0.090723D-06 0.0
8 0 0.049118D-06 0.0
2 1 −0.000187D-06 0.001195D-06
3 1 2.030137D-06 0.248131D-06
4 1 −0.536244D-06 −0.473772D-06
5 1 −0.062727D-06 −0.094195D-06
6 1 −0.076104D-06 0.026900D-06
7 1 0.280287D-06 0.094777D-06
8 1 0.023334D-06 0.058499D-06
2 2 2.439261D-06 −1.400266D-06
3 2 0.904706D-06 −0.618923D-06
4 2 0.350670D-06 0.662571D-06
5 2 0.652459D-06 −0.323334D-06
6 2 0.048328D-06 −0.373816D-06
7 2 0.329760D-06 0.093194D-06
8 2 0.080071D-06 0.065519D-06
3 3 0.721145D-06 1.414204D-06
4 3 0.990869D-06 −0.200987D-06
5 3 −0.451837D-06 −0.214954D-06
6 3 0.057021D-06 0.008890D-06
7 3 0.250502D-06 −0.217320D-06
8 3 −0.019252D-06 −0.086286D-06
4 4 −0.188481D-06 0.308848D-06
5 4 −0.295123D-06 0.049741D-06
6 4 −0.086228D-06 −0.471405D-06
7 4 −0.275541D-06 −0.124142D-06
8 4 −0.244358D-06 0.069857D-06
(continued)
(continued)
n m C nm Snm
5 5 0.174832D-06 −0.669393D-06
6 5 −0.267112D-06 −0.536410D-06
7 5 0.001644D-06 0.018075D-06
8 5 −0.025498D-06 0.089090D-06
6 6 0.009502D-06 −0.237262D-06
7 6 −0.358843D-06 0.151778D-06
8 6 −0.065859D-06 0.308921D-06
7 7 0.001380D-06 0.024129D-06
8 7 0.067263D-06 0.074813D-06
8 8 −0.123971D-06 0.120441D-06
The relations between the normalised coefficients C nm and Snm and their
non-normalised counterparts Cnm and Snm are given below (from Lambeck [17])
C nm Snm
Cnm ¼ Snm ¼
N nm N nm
where the normalising factor Nnm is

12
ðn mÞ!
N nm ¼ ð2 dm0 Þð2n þ 1Þ
ðn þ mÞ!
In the preceding expression, dm0 (the Kronecker delta function) is equal to unity
if m = 0 and is equal to zero otherwise.
Likewise, the normalised associated Legendre function of the first kind Pnm
(over-lined) is used in place of Pnm.
The relation between Pnm and Pnm is given below (from Lambeck [17])
Pnm ¼ Pnm N nm
By so doing, the squared norm of a spherical harmonic Pnm (sin u*) eimk over a
unit sphere is 4p.
Using the normalised Stokes coefficients and the normalised associated
Legendre function of the first kind, the acceleration vector f r″ = grad(V), due to
the actual potential of gravitation V(r, k, u*) of the Earth, can be written as follows
( )
l X
1 X n n
rE
00
r ¼ grad C nm cosðmkÞ þ Snm sinðmkÞ Pnm ðsinu Þ
r n¼0 m¼0 r
where r″ d2r/dt2 is the second derivative of the position vector r with respect to
time. Reference [23] gives the values of the normalised Stokes coefficients C nm and
Snm , for the EGM96 gravity model, truncated at n = m = 18.
There are recurrence relations that make it easy to compute the associated
Legendre function of the first kind Pn,m(sin u*) starting from those of lower degree
n and order m.
According to Losch and Seufer (24), such recurrence relations are:
Pn þ 1;0 ðsin u Þ ¼ ð2n þ 1Þ sin u Pn;0 ðsin u Þ nPn1;0 ðsin u Þ

Pn;n ðsin u Þ ¼ ð2n 1Þ cos u Pn1;n1 ðsin u Þ
Pn;m ðsin u Þ ¼ Pn2;m ðsin u Þ þ ð2n 1Þ cos u Pn1;m1 ðsin u Þ
with the following starting values
P0,0 = 1
P1,0 = sin u* P1,1 = cos u*
P2,0 = ½(3 sin2u* − 1) P2,1 = 3 cos u* sin u* P2,2 = 3 cos2u*
Losch and Seufer point out that the formulae given above become numerically
unstable for n and m greater than 120.
In order to avoid instability, Montenbruck and Gill [2] propose to proceed as
follows. Starting from P0,0 = 1, all polynomials Pm,m(sin u*) up to the desired
degree and order are first computed from
Pm;m ðsin u Þ ¼ ð2m 1Þ cos u Pm1;m1 ðsin u Þ
With these results, the remaining values are computed from
Pm þ 1;m ðsin u Þ ¼ ð2m þ 1Þ sin u Pm;m ðsin u Þ
and from the following recursion
ð2n 1Þðsin u ÞPn1;m ðsin u Þ ðn þ m 1ÞPn2;m ðsin u Þ

Pn;m ðsin u Þ ¼
nm
for n > m + 1. To the same end, Press et al. (25) propose to compute the two
starting values Pm,m(sin u*) and Pm+1,m(sin u*) as follows
m
Pm;m ðsin u Þ ¼ ð2m 1Þ!! 1 sin2 u 2
Pm þ 1;m ðsin u Þ ¼ ð2m þ 1Þðsin u ÞPm;m ðsin u Þ
where the notation k!! (called double factorial) indicates the product of all odd
integers less than or equal to k. In terms of the ordinary factorial, the double
factorial is given by
ð2mÞ!
ð2m 1Þ!! ¼
m!2m
When the two starting values Pm,m(sin u*) and Pm+1,m(sin u*) are known, Press
et al. compute Pn,m(sin u*) in the same way as that indicated by Montenbruck and
Gill [2], that is, as follows
ðn mÞPn;m ðsin u Þ ¼ ð2n 1Þðsin u ÞPn1;m ðsin u Þ

ðn þ m 1ÞPn2;m ðsin u Þ
Press et al. [26] also give a c++ function, called plgndr, which computes the
associated Legendre function of the first kind Pn,m(sin u*) for given values of n, m,
and sin u*.
The recurrence relations that can be used to compute the normalised associated
Legendre function of the first kind Pn;m (sin u*), starting from those of lower degree
n and order m, are those suggested by Holmes and Featherstone [27] and also by
Losch and Seufer [27]. They are given below. First, the sectorial (n = m) nor-
malised functions Pm,m(sin u*) are computed by means of

1
2m þ 1 2
Pm;m ðsin u Þ ¼ ðcos u Þ Pm1;m1 ðsin u Þ
2m
for all m > 1, starting from P0;0 (sin u*) = 1 and P1;1 (sin u*) = (3)½ cos u*.
Then, after computing Pm;m (sin u*), the non-sectorial (n > m) normalised
functions Pn;m (sin u*) are computed by means of
Pn;m ðsin u Þ ¼ ðsin u Þan;m Pn1;m ðsin u Þ bn;m Pn2;m ðsin u Þ
where
1 1
ð2n 1Þð2n þ 1Þ 2 ð2n þ 1Þðn þ m 1Þðn m 1Þ 2
an;m ¼ bn;m ¼
ðn mÞðn þ mÞ ðn mÞðn þ mÞð2n 3Þ
for all n > m. The recurrence relations shown above for the associated Legendre
functions can be used together with the well-known addition formulae
cos½ðm þ 1Þk ¼ cosðmkÞ cos k sinðmkÞ sin k

sin½ðm þ 1Þk ¼ sinðmkÞ cos k þ cosðmkÞ sin k
for the sine and cosine functions whose argument depends on longitude (k). This
makes it possible to write a single recurrence relation for the geopotential, as will be
shown below. Since l/r = (l/rE)(rE/r) is an identity, then Cunningham (26) and
Montenbruck and Gill [2] define the following quantities
r n þ 1
½Pnm ðsin u Þ cosðmkÞ
E
V nm ¼
r
r n þ 1
½Pnm ðsin u Þ sinðmkÞ
E
W nm ¼
r
and express the potential of gravitation V of the Earth as follows

X
l 1 X n
Vðr; k; u Þ ¼ ðCnm V nm þ Snm W nm Þ
rE n¼0 m¼0
where rE is the equatorial radius and l is the gravitational parameter of the Earth.
Let x, y, and z be the Cartesian co-ordinates of a point P in the geocentric-equatorial
Earth-fixed (that is, non-inertial) reference system. These co-ordinates are related to
the spherical co-ordinates r, k, and u* of P as follows
x ¼ r cos u cos k
y ¼ r cos u sin k
z ¼ r sin u
Following Montenbruck and Gill [2], the quantities Vnm and Wnm satisfy the
following recurrence relations
xr yrE
E
V mm ¼ ð2m 1Þ V m1;m1 W m1;m1
r2
xr r2
E yrE
W mm ¼ ð2m 1Þ 2 W m1;m1 2 V m1;m1
r r
and
2n 1 zrE n þ m 1 rE 2
V nm ¼ V n1;m V n2;m
n m r2 nm r

2n 1 zrE
n þ m 1 rE 2
W nm ¼ W n1;m W n2;m
nm r 2 nm r
which result from the recurrence relations for the associated Legendre functions and
also from the trigonometric formulae for addition of angles.
The second pair of relations for Vnm and Wnm also holds for n = m + 1, if Vm−1,m
and Wm−1,m are set to zero.
The starting values for these recurrence relations are known. They are
rE
V 00 ¼ W 00 ¼ 0
r
Following Holmes and Featherstone [27], we represent the scheme for com-
puting the successive pairs of quantities Vnm and Wnm by means of a lower trian-
gular matrix, shown in the following figure, where each circle corresponds to a
particular combination of degree (n) and order (m).
For the purpose of computing the quantities Vnm and Wnm (0 m n) up to

some chosen nmax, the zonal terms Vn0 are computed first, by means of
n m r2 nm r
nm r 2 nm r
with m = 0 and n = 1, 2, …, nmax. There is no need to compute the zonal terms

Wn0, because they are, all of them, equal to zero. By so doing, it is possible to fill
the leftmost (corresponding to m = 0) column of the matrix shown above.
Now, we start again from
rE
V 00 ¼ W 00 ¼ 0
r
and compute the pair (V11, W11) by using the recurrence relations
xr yrE
E
r2
xr r2
E yrE
W mm ¼ ð2m 1Þ 2 W m1;m1 2 V m1;m1
r r
with m = 1. When the pair (V11, W11) is known, we fill the column corresponding to
m = 1 by using again the recurrence relations
n m r2 nm r
nm r 2 nm r
with m = 1 and n = 2, 3, …, nmax. Now, starting again from the pair (V11, W11), we
compute the pair (V22, W22) by means of
xr yrE
E
r2
xr r2
E yrE
W mm ¼ ð2m 1Þ 2 W m1;m1 2 V m1;m1
r r
with m = 2. The triangular matrix indicated above is filled column by column until
the degree (n) reaches the chosen value nmax. The expressions written above make
it possible to compute the components x″, y″, and z″ of the acceleration vector
r″ = grad(V) due to the Earth gravitational field, as follows
X X X
x00 ¼ x00nm y00 ¼ y00nm z00 ¼ z00nm
n;m n;m n;m
where the partial accelerations are given below

l
x00nm ¼ Cn0 V n þ 1;1 ðm ¼ 0Þ
rE2
l1 ðn m þ 2Þ!
x00nm ¼ 2 ½ Cnm V n þ 1;m þ 1 Snm W n þ 1;m þ 1 þ
rE 2 ðn mÞ!

þ C nm V n þ 1;m1 þ Snm W n þ 1;m1 ðm [ 0Þ
l
y00nm ¼ 2
C n0 W n þ 1;1 ðm ¼ 0Þ
rE
l1 ðn m þ 2Þ!
y00nm ¼ 2 ½ C nm W n þ 1;m þ 1 þ Snm V n þ 1;m þ 1 þ
rE 2 ðn mÞ!

C nm W n þ 1;m1 þ Snm V n þ 1;m1 ðm [ 0Þ
l
z00nm ¼ 2
ðn m þ 1Þ C nm V n þ 1;m Snm W n þ 1;m
rE
in accordance with Ref. [2].

Montenbruck and Gill [2] note that the quantities Vnm and Wnm are required up to
degree n + 1 and order m + 1, in order for the partial accelerations due to
geopotential coefficients up to Cnn and Snn to be calculated.
The expressions given above for x″, y″, and z″ make it possible to compute the
acceleration vector frot r″rot, due to the gravitational attraction of the Earth, in a
non-inertial geocentric reference system xyz which rotates with the Earth.
When a high accuracy is not required, it is not necessary to integrate numerically
the Lagrange planetary equations for the purpose of computing the changes in the
right ascension of the ascending node (X) and in the argument of perigee (x)
induced by the non-spherical shape of the Earth, as the following example will
show. Let us consider an artificial satellite orbiting around the Earth. At a given
epoch t0, that is, on the 6th of November 2000, 22h:34m:30s UT1, let
r0 ¼ 3525:5 uX 4313:7 uY þ 4651:7 uZ

v0 ¼ 4:0755 uX þ 2:6425 uY þ 5:5324 uZ
be the position and velocity vectors of the satellite in geocentric-equatorial

co-ordinates X0, Y0, Z0, X′0, Y′0, and Z′0 expressed in units of, respectively, km and
km/s. We want to compute the position and velocity vectors of the same satellite at
an epoch tf five days (or 5 86400 s) later than t0, in the same reference system,
under the following hypotheses: (1) no other perturbations than that due to the Earth
oblateness affect the orbit; (2) the oblate ellipsoid model is a valid approximation to
the actual gravitational field of the Earth; and (3) the first-order secular (and not the
periodic) effects of the non-spherical Earth are taken into account. In these con-
ditions, as has been shown above, the major semi-axis (a), the eccentricity (e), and
the inclination (i) of the orbit are not perturbed. First of all, we compute the orbital
elements of the satellite at epoch t0, as follows.
The major semi-axis (a0) at epoch t0 can be computed by means of the vis-viva
integral, as follows
v20 2 1
¼
lE r0 a0
where lE = 3.986 105 km3/s2 is the gravitational parameter of the Earth. In the
present case, r0 and v20 are respectively
h i12
r0 ¼ 3525:52 þ ð4313:7Þ2 þ 4651:72 ¼ 7257:8 km
v20 ¼ ð4:0755Þ2 þ 2:64252 þ 5:53242 ¼ 54:200 km2 =s2
Since the major semi-axis, the eccentricity, and the inclination of the orbit are
not perturbed, then the subscript 0 is not strictly necessary with these quantities and
may also be omitted. The vis-viva integral, solved for the major semi-axis, yields
1 1
a0 ¼ ¼ ¼ 7163:8 km
2 v20 2 54:200

r0 lE 7257:8 398600
Remembering that the moment-of-momentum vector per unit mass

h0 ¼ r 0 v 0
also remains constant with time, there results

h0 ¼ hX0 uX þ hY0 uY þ hZ0 uZ ¼ Y0 Z00 Z0 Y00 uX þ Z0 X00 X0 Z00 uY

þ X0 Y00 Y0 X00 uZ ¼ ð4313:7 5:5324 4651:7 2:6425ÞuX
þ ð4651:7 4:0755 3525:5 5:5324ÞuY þ ð3525:5 2:6425
4313:7 4:0755ÞuZ ¼ 36157uX 38462uY 8264:4uZ
The magnitude of h0 is
1
h i12
h0 ¼ ðh0 h0 Þ2 ¼ ð36157Þ2 þ ð38462Þ2 þ ð8264:4Þ2 ¼ 53432 km2 =s
The (constant) semi-latus rectum results from
h20 534322
p0 ¼ ¼ ¼ 7162:6 km
lE 398600
The (constant) scalar eccentricity results from

p0 ¼ a0 1 e20
which, solved for e0, yields

1
1
p0 2 7162:6 2
e0 ¼ 1 ¼ 1 ¼ 0:013130
a0 7163:8
The following expression (derived in Chap. 1, Sect. 1.1)

v0 h0 r0
e0 ¼
lE r0
makes it possible to compute the eccentricity vector (e0), as follows
ðv0 h0 Þ=lE ¼ ½ð2:6425 8264:4 þ 5:5324 38462ÞuX

þ ð5:5324 36157 4:0755 8264:4ÞuY

þ ½ð4:0755 38462 þ 2:6425 36157ÞuZ = 3:986 105
¼ 0:47905 uX 0:58635 uY þ 0:63296 uZ
r0 =r0 ¼ ð3525:5 uX 4313:7 uY þ 4651:7 uZ Þ=7257:8

¼ 0:48575 uX 0:59435 uY þ 0:64093 uZ
Thus, the eccentricity vector results from
e0 ¼ ðv0 h0 Þ=lE r0 =r0 ¼ ð0:47905 0:48575Þ uX þ ð0:58635 þ 0:59435Þ uY

þ ð0:63296 0:64093Þ uZ ¼ 0:0067 uX þ 0:008 uY 0:00797 uZ
As a check, the scalar eccentricity may be computed again by evaluating the

magnitude of the eccentricity vector. This yields
1
h i12
e0 ¼ ðe0 e0 Þ2 ¼ ð0:0067Þ2 þ 0:0082 þ ð0:00797Þ2 ¼ 0:01313
The orbital plane is constantly inclined with respect to the equatorial plane by

h0Z 8264:4
i0 ¼ arccos ¼ arccos ¼ 98 :898 ¼ 1:7261 radians
h0 53432
The right ascension of the ascending node at epoch (X0) results from computing
first the node vector (n0) defined by
n0 ¼ uZ h0
Thus there results
n0 ¼ uZ h0 ¼ n0X uX þ n0Y uY ¼ h0Y uX þ h0X uY ¼ ð38462Þ uX þ ð36157Þ uY

¼ 38462 uX 36157 uY
The magnitude of the node vector is
1 h i12
n0 ¼ n20X þ n20Y 2 ¼ 384622 þ ð36157Þ2 ¼ 52789
In the general case, the angle X0 results from

n0X
X0 ¼ arccos ðn0Y
0Þ
n0

n0X
X0 ¼ 360 arccos ðn0Y \0Þ
n0
In the present case, n0Y = −36157 < 0; hence

n0X 38462
X0 ¼ 360 arccos ¼ 360 arccos ¼ 316 :77
n0 52789
In the general case, the argument of perigee at epoch (x0) results from

n0 e0
x0 ¼ arccos ðe0Z
0Þ
n0 e 0

n 0 e0
x0 ¼ 360 arccos ðe0Z \0Þ
n0 e 0
In the present case, e0Z = −0.00797 < 0; hence

38462 0:0067 36157 0:008
x0 ¼ 360 arccos ¼ 217 :90
52789 0:01313
In the general case, the true anomaly at epoch (/0) results from

e0 r 0
/0 ¼ arccos ðr0 v0
0Þ
e0 r 0

e0 r 0
/0 ¼ 360 arccos ðr0 v0 \0Þ
e0 r0
r0 v0 ¼ 3525:5 ð4:0755Þ þ ð4313:7Þ 2:6425 þ 4651:7 5:5324

¼ 32:062\0
Therefore, the true anomaly at epoch is

0:0067 3525:5 0:008 4313:7 0:00797 4651:7
/0 ¼ 360 arccos
0:013130 7257:8
¼ 182 :53
The orbital period is

12
12
a3 7163:83
T0 ¼ 2p 0 ¼ 2 3:1416 ¼ 6034:3 s
lE 398600
The mean motion is

2p 2 3:1416
n0 ¼ ¼ ¼ 0:0010412 radians/s
T0 6034:3
The eccentric anomaly at epoch (Æ0) is computed by using the following for-
mula given in Chap. 1, Sect. 1.3:

1

1

Æ0 1 e0 2 /0 1 0:01313 2 182 :53
tan ¼ tan ¼ tan ¼ 44:695
2 1 þ e0 2 1 þ 0:01313 2
Hence
Æ 0 ¼ 2 arctanð44:695Þ ¼ 182 :56 ¼ 3:1863 radians
Since the mean motion (n0) and the eccentric anomaly (Æ0) are known, Kepler’s
equation can be used to compute the time (t1) lapsed from the moment of passage at
perigee to the initial epoch (t0), as follows
n0 t1 ¼ Æ 0 e0 sin Æ 0
The preceding equation, solved for t1, yields
Æ 0 e0 sin Æ 0 3:1863 0:01313 sin 3:1863

t1 ¼ ¼ ¼ 3060:7 s
n0 0:0010412
The final epoch (tf) is computed by adding five days (5 86400 s) to t1. This
yields
tf ¼ 3060:7 þ 5 86400 ¼ 4:3506 105 s
The number of periods (nT) since the satellite passed at perigee in its first orbit is
tf 4:3506 105
nT ¼ ¼ ¼ 72:098
T0 6034:3
In other words, the satellite has completed 72 orbits at the final epoch tf. It has
travelled along a fraction of orbit at the initial epoch t0. An orbit begins at the
moment of a passage at perigee and ends at the moment of the next passage at
perigee. Let t73 denote the time spent by the satellite from the beginning of its 73rd
orbit to the final epoch tf. This time results from
t73 ¼ ð72:098 72:000ÞT0 ¼ 0:098 6034:3 ¼ 590:44 s
The mean anomaly corresponding to this time in the 73rd orbit is
M73 ¼ n0 t73 ¼ 0:0010412 590:44 ¼ 0:61479 radians
The eccentric anomaly Æ73 corresponding to M73 results from Kepler’s equation
M73 ¼ Æ 73 e0 sin Æ 73
This equation is solved iteratively for Æ73, as will be shown below. For con-
venience, we set Æ73 x and search a zero of the following function
f ðxÞ ¼ x 0:013130 sinðxÞ 0:61479
A good approximation to the unknown value of Æ73 x is computed below by

means of the following formula given by Battin [11]:
e sin M
Æ Mþ
1 sinðM þ eÞ þ sin M
In the present case, we have
0:01313 sin 0:61479

Æ 73 0:61479 þ ¼ 0:62634 radians
1 sinð0:61479 þ 0:01313Þ þ sin 0:61479
Therefore, we search the unknown value of Æ73 x within the following

interval
0:62 x 0:63
with an accuracy of five significant figures. At the lower end (0.62) and at the upper
end (0.63) of the interval, there results
f ð0:62Þ ¼ 0:62 0:013130 sinð0:62Þ 0:61479 ¼ 0:0024190

f ð0:63Þ ¼ 0:63 0:013130 sinð0:63Þ 0:61479 ¼ 0:0074745
Since the product f(0.62)f(0.63) is less than zero, then there is at least a zero of
the function f(x) within the interval 0.62 x 0.63. We choose arbitrarily the
mid-point (x = 0.625) as another point of this interval and compute
f ð0:625Þ ¼ 0:625 0:013130 sinð0:625Þ 0:61479 ¼ 0:0025277
Then we set
x2 ¼ 0:620 f2 f ðx2 Þ ¼ 0:0024190

x0 ¼ 0:625 f0 f ðx0 Þ ¼ 0:0025277
x1 ¼ 0:630 f1 f ðx1 Þ ¼ 0:0074745
and
h1 ¼ x1 x0 ¼ 0:630 0:625 ¼ 0:005

h2 ¼ x0 x2 ¼ 0:625 0:620 ¼ 0:005
c ¼ h2 =h1 ¼ 0:005=0:005 ¼ 1
Now we compute the coefficients
A ¼ ½cf1 f0 ð1 þ cÞ þ f2 =½ch21 ð1 þ cÞ ¼ ½1 0:0074745 0:0025277 ð1 þ 1Þ

þ ð 0:0024190Þ=½1 0:0052 ð1 þ 1Þ ¼ 0:002

B ¼ f1 f0 Ah21 =h1 ¼ ð0:0074745 0:0025277 0:002 0:0052 Þ=0:005
¼ 0:98935
C ¼ f0 ¼ 0:0025277
of the interpolating parabola f(x) = A(x − x0)2 + B(x − x0) + C and evaluate the
estimated root of f(x) = 0 as follows
2C
x ¼ x0 1
B ðB2 4AC Þ2
where the sign plus or minus is chosen so that the denominator should have the
maximum absolute value (that is, if B > 0, choose plus; if B < 0, choose minus; if
B = 0, choose either). In the present case (B > 0), we chose plus, and therefore
2 0:0025277
x ¼ 0:625 1 ¼ 0:62245
0:98935 þ ð0:989352 4 0:002 0:0025277Þ2
We use this value to compute
f ðxÞ ¼ f ð0:62245Þ ¼ 0:62245 0:013130 sinð0:62245Þ 0:61479 ¼ 0:00000485
Since we search the unknown value of x with an accuracy of five significant

figures, then we accept Æ73 = 0.62245 radians = 35°.664 as the correct solution of
Kepler’s equation.
The true anomaly corresponding to Æ73 = 0.62245 radians results from

12

1

/ 1 þ e0 Æ 73 1 þ 0:01313 2 0:62245
tan 73 ¼ tan ¼ tan ¼ 0:32593
2 1 e0 2 1 0:01313 2
which, solved for /73, yields
/73 ¼ 0:63015 radians ¼ 36 :105
Since we consider only the first-order secular (that is, non-periodic) effects, then
the perturbation due to the ellipsoidal Earth does not affect the major semi-axis (a0),
the eccentricity (e0), and the inclination (i0) of the orbit of a satellite. Therefore, the
radius vector (r73) corresponding to the true anomaly /73 can be computed by using
the constant values of a0, e0 and i0, as follows
p0 7162:6
r73 ¼ ¼ ¼ 7087:4 km
1 þ e0 cos /73 1 þ 0:01313 cos 36 :105
For the purpose of expressing the position and velocity vectors in the perifocal
system of reference Oxyz, we use the formulae of Chap. 1, Sect. 1.9

12
lE
v¼ ð sin /Þux þ ðe þ cos /Þuy
p
Since
r73 cos /73 ¼ 7087:4 cos 36 :105 ¼ 5726:2 km

r73 sin /73 ¼ 7087:4 sin 36 :105 ¼ 4176:3 km
then the position vector of the satellite in the perifocal system of reference,
expressed in km, at the final epoch, is
r73 ¼ ðr 73 cos /73 Þux þ ðr 73 sin /73 Þuy ¼ 5726:2ux þ 4176:3uy
Likewise, since

12
1
lE 398600 2
ð sin /73 Þ ¼ ð sin 36 :105Þ ¼ 4:3959 km/s
p0 7162:6

12
1
lE 398600 2
ðe0 þ cos /73 Þ ¼ ð0:01313 þ cos 36 :105Þ ¼ 6:1251 km/s
p0 7162:6
then the velocity vector of the satellite in the perifocal system of reference,
expressed in km/s, at the final epoch, is
v73 ¼ 4:3959ux þ 6:1251uy
Since the true anomaly of the satellite has been updated, it remains to update the
right ascension of the ascending node and the argument of perigee.
The rate of change of the right ascension of the ascending node with time has
been shown above to be
3 r 2 cos i
X0 ¼ nC20
E
2 a ð1 e2 Þ2
where n is the mean motion of the satellite, rE is the equatorial radius of the Earth,
and C20 is the principal zonal coefficient. In the present case, there results

2
6378:1 cos 1:7261
X0 ¼ 1:5 0:0010412 ð0:0010826Þ
7163:8 ð1 0:013132 Þ2
¼ 2:0739 107 radians=s ¼ 1:1882 105 degrees=s
Therefore, the right ascension of the ascending node at the final epoch is
X73 ¼ X0 þ X0 Dt ¼ 316 :77 þ 1:1882 105 5 86400 ¼ 321 :90
The rate of change of the argument of perigee with time has been shown to be
3 r 2 1 5 cos2 i
x0 ¼ nC 20
E
4 a ð 1 e2 Þ 2

2
0 6378:1 1 5 cos2 1:7261
x ¼ 0:75 0:0010412 ð0:0010826Þ
7163:8 ð1 0:013132 Þ2
¼ 5:9023 107 radians=s ¼ 3:3817 105 degrees=s
Therefore, the argument of perigee at the final epoch is
x73 ¼ x0 þ x0 Dt ¼ 217 :90 3:3817 105 5 86400 ¼ 203 :29
Now, for convenience, we change the notation as follows: r73 = r, v73 = v,

X73 = X, x73 = x, and i0 = i.
We use the following rotation matrix R shown in Chap. 1, Sect. 1.9:



The components X, Y, and Z of the position vector r of the satellite, in the

geocentric-equatorial reference system, at the final epoch, result from
pre-multiplying the column vector whose components are x, y, and 0 (the com-
ponents of the position vector of the satellite in the perifocal system, at the same
epoch) by the rotation matrix R, as follows
X ¼ ðcos X cos x sin X sin x cos iÞx þ ð cos X sin x

sin X cos x cos iÞy ¼ ðcos 321 :90 cos 203 :29 sin 321 :90
sin 203 :29 cos 98 :898Þ 5726:2 þ ð cos 321 :90 sin 203 :29
sin 321 :90 cos 203 :29 cos 98 :898Þ 4176:3 ¼ 2257:5 km
Y ¼ ðsin X cos x þ cos X sin x cos iÞx þ ð sin X sin x

þ cos X cos x cos iÞy ¼ ðsin 321 :90 cos 203 :29 þ cos 321 :90
sin 203 :29 cos 98 :898Þ 5726:2 þ ð sin 321 :90 sin 203 :29
þ cos 321 :90 cos 203 :29 cos 98 :898Þ 4176:3 ¼ 2968:9 km
Z ¼ ðsin x sin iÞx þ ðcos x sin iÞy ¼ ðsin 203 :29 sin 98 :898Þ 5726:2
þ ðcos 203 :29 sin 98 :898Þ 4176:3 ¼ 6026:7 km
Likewise, the components vX, vY, and vZ (or X′, Y′, and Z′) of the velocity vector
v of the satellite, in the geocentric-equatorial reference system, at the final epoch,
result from pre-multiplying the column vector whose components are vx, vy, and 0
(the components of the velocity vector of the satellite in the perifocal system, at the
same epoch) by the rotation matrix R, as follows
vX ¼ ðcos X cos x sin X sin x cos iÞvx þ ð cos X sin x

sin X cos x cos iÞvy ¼ ðcos 321 :90 cos 203 :29 sin 321 :90
sin 203 :29 cos 98 :898Þ ð4:3959Þ þ ð cos 321 :90 sin 203 :29
sin 321 :90 cos 203 :29 cos 98 :898Þ 6:1251 ¼ 5:4544 km=s
vY ¼ ðsin X cos x þ cos X sin x cos iÞvx þ ð sin X sin x

þ cos X cos x cos iÞvy ¼ ðsin 321 :90 cos 203 :29 þ cos 321 :90
sin 203 :29 cos 98 :898Þ ð4:3959Þ þ ð sin 321 :90 sin 203 :29
þ cos 321 :90 cos 203 :29 cos 98 :898Þ 6:1251 ¼ 3:5122 km=s
vZ ¼ ðsin x sin iÞvx þ ðcos x sin iÞvy ¼ ðsin 203 :29 sin 98 :898Þ ð4:3959Þ
þ ðcos 203 :29 sin 98 :898Þ 6:1251 ¼ 3:8412 km=s
The required position and velocity vectors of the satellite, in the

geocentric-equatorial reference system, at the final epoch tf, are then
rf ¼ 2257:5uX þ 2968:9uY 6026:7uZ

vf ¼ 5:4544uX 3:5122uY 3:8412uZ
3.3 The Changes of Orientation of the Earth Axis
As has been shown in Chap. 2, Sect. 2.4, the angular velocity at which the Earth
rotates about its axis is not constant in time, nor is the axis of rotation of the Earth
fixed in space. In addition, owing to tectonic movements, which are large-scale
3.3 The Changes of Orientation of the Earth Axis 357
motions of the Earth lithosphere, neither the shape nor the relative positions of
locations on the Earth surface are fixed.
Let us consider the changes of orientation of the Earth rotation axis. As a result
of these changes, the astronomers must refer the angles (declination and right
ascension) of the stars to some epochs, the most common of which are 31
December 1949 at 22:09 UT (B1950.0) and 1 January 2000 at 12:00 TT (J2000.0).
The zero of right ascension is the direction of the vernal equinox (first point of
Aries), that is, the direction of the line where the Earth equatorial plane intersects
the ecliptic plane. Since the direction of the Earth rotation axis varies in time with
respect to the stars, so does the direction of the vernal equinox. The variation with
time of the direction of the Earth rotation axis depends on the following four
components: precession, nutation, celestial pole offset, and polar motion. These
components are considered in the following sections.
Neither of the two fundamental planes (the ecliptic and the equatorial plane) is
fixed with respect to the distant stars. By precession, we mean the conical motion of
the axis of rotation of the Earth about the pole of the ecliptic plane, as shown in the
following figure, which is due to the courtesy of Wikimedia [28]. This motion is
due to torques exerted on the rotating oblate Earth by the Moon and the Sun. These
torques induce a slow (each full cycle takes about 26000 years) and continuous
change in the orientation of the Earth rotation axis along a pair of cones, joined at
their vertices, whose axis is perpendicular to the plane of the ecliptic.
The direction in the sky to which the Earth rotation axis points through each
26000-year cycle describes a big circle, whose radius subtends an angle of 23°.5.
As a result of this motion, the Earth wobbles slowly like a top, causing the
position of the vernal equinox to shift westward (hence, the name of precession)
around the sky. The rate of this shift is a little more than 50 arcseconds per year, as
will be shown below. At the present time, the vernal equinox is in the constellation
Pisces and moves towards Aquarius.
For the same reason, the star which is nearest to the north celestial pole changes
with time. Consequently, Polaris will not always be the North Star. In about
14000 years, the North Star will again be Vega (also known as a Lyrae), which was
the North Star about 12000 years ago. After a lapse of 26000 years from the present
time, the Earth rotation axis will again point towards Polaris. This fact is shown in
the preceding figure, due to the courtesy of Wikimedia [29].
In this figure, two constellations are shown, namely Ursa Minor (whose a star,
Polaris, is the North Star now at 2000 AD) and Draco (whose a star, Thuban, was
the North Star at about 3000 BC).
In about 129 BC, according to Heath [30], Hipparchus of Nicaea estimated the
precession of the Earth rotation axis by measuring the ecliptic longitude of the star
Spica (also known as a Virginis) during lunar eclipses.
His observations showed that Spica’s ecliptic longitude was about 6° west of the
autumnal equinox. Then, by comparing these observations with those which had
previously been performed by Timocharis of Alexandria, Hipparchus found that the
same angular distance had been made 8° by Timocharis. Since the motion
amounted to 2° in the period between Timocharis’s observations, performed in 283
or 295 BC, and 129 or 128 BC, that is, in a period of 154 or 166 years, then the
estimated motion of precession was
3600 3600
2 ¼ 46:7500 =year or 2 ¼ 43:1100 =year
283 129 295 128
These values are compared with the value 5028.796195″/Julian century at

J2000.0 [31]. By using this value, the precession period of the Earth results
3600
T ¼ 360 25722 years
50:288
In addition to the long-period motion (precession) described above, the Earth

rotation axis has small short-period oscillations, superimposed to the precession,
with respect to the distant stars. This motion is known as nutation and is shown in
the following figure, due to the courtesy of Wikimedia [32]. This is because the
orbital plane of the Moon is inclined 5°.1 with respect to the ecliptic plane. The
straight line marking the intersection of these planes is the line of nodes of the
Moon. Following Huggett [33], this line regresses (with respect to the sense of
rotation of the Earth and revolution of the Moon around the Earth), taking on
average 18.6134 years to complete one cycle. This is the lunar nodal (or Metonic)
cycle. In addition, the longitude of the lunar perigee progresses (that is, revolves
eastward) in the orbital plane of the Moon, taking on average 8.849 years to
complete one cycle (apsides cycle of the Moon). The nodal cycle of the Moon
causes small variations in the lunar torque as the perigee and apogee come into the
plane of the axial tilt of the Earth. These variations are sufficient to make the Earth
rotation axis nod up and down a little.
In the preceding figure, R is the rotation axis of the Earth, P is the precession
motion, and N is the nutation motion. The nutations are small oscillations on the
surface of the cone of precession. There is more than one period in the nutation
motion. The primary nutation cycle imposes a 9″.18 increase or decrease of the
Earth rotation axis over an 18.6-year cycle. Secondary nutations induce a change in
inclination amounting to less than 1″ over a cycle of 18.6/2 = 9.3 years. In addi-
tion, there are annual and semi-annual nutation cycles (caused by variations of the
solar couple at perihelion and aphelion) and 13.6-day cycles (caused by variations
of the solar couple at perigee and apogee), as has been shown in Ref. [33].
Finally, internal motions and consequent deformations of the Earth cause small
variations in the mutual distances between observation points located on the Earth
surface. Consequently, the Earth rotation axis moves with respect to the Earth
surface. This motion is called polar motion and is shown in the following figure
(courtesy of the IERS [34], relating to the polar motion in the time interval 1
November 2010–10 November 2011. In this figure, mas stands for milliarcseconds
(on the Earth surface, 100 milliarcseconds are nearly 3 m).
The principal component of the polar motion is the so-called Chandler wobble
(after S.C. Chandler, who first observed it in 1891), which is a nearly circular
motion of the Earth pole around the celestial pole.
According to Gross [35], polar motion consists largely of:

(1) a forced annual wobble having a nearly constant amplitude of about 100
milliarcseconds;
(2) the free Chandler wobble having a variable amplitude ranging between 100 and
200 milliarcseconds;
(3) quasi-periodic variations on decadal timescales having amplitudes of about 30
milliarcseconds (known as the Markowitz wobble);
(4) a linear trend having a rate of about 3.5 milliarcseconds per year; and
(5) smaller amplitude variations occurring on all measurable timescales.
In particular, the Chandler wobble has a period of 433 days and is excited by a
combination of atmospheric and oceanic processes, in which ocean-bottom pressure
fluctuations are the principal cause of excitation [35].
3.4 The Change of Co-ordinates Due to Precession
The purpose of the present section is to show how to determine the precession
matrix P(t) which takes account of the motion of precession of the Earth rotation
axis from J2000.0 to another date t. This motion produces a change in the ecliptic
longitude of an object. The gravitational attraction exerted by the Sun and Moon
changes the direction of the rotation axis of the Earth, but has no effect on the
ecliptic plane. However, the latter is subject to change with time as a result of forces
exerted by the planets, as will be shown in the following sections.
In accordance with Vallado [36], when the precession matrix P(t) is known, it is
possible to convert a position vector in the Earth Centred Inertial J2000 system
(defined on the basis of the mean equinox of J2000) to the corresponding vector in
the Earth Centred Inertial mean-of-date system (defined by the mean equinox of the
date of interest).
Let rJ2000mod be a vector whose components are known in the Earth Centred
Inertial J2000-mean-equinox system. Let rECImod be the same vector whose com-
ponents are to be determined in the Earth Centred Inertial mean-equinox-of-date
system. The components of rECImod result from those of rJ2000mod by means of the
following expression
rECImod ¼ PðtÞ rJ2000mod
The precession matrix P(t) contains nine quantities pij depending on time. Let
JDE be the Julian ephemeris day of interest, computed as shown in Chap. 2,
Sect. 2.4. Let
t ¼ JDE 2451545:0
JDE 2451545:0
T¼
36525
be the time (expressed in, respectively, Julian days and Julian centuries) elapsed
between the Julian ephemeris day of interest and J2000.0. JDE is the Julian day
(JD) corresponding to an instant of time measured in the Terrestrial Timescale. For
convenience of the reader, we give below the principal formulae of Chap. 2,
Sect. 2.4. First, JD results from
UT1
JD ¼ JDN þ
24
where JDN (Julian Day Number) is given by

n h m io
JDN ¼ 367y INT 1:75 y þ INT þ 0:75

12
275
þ INT m þ d þ 1721013:5
9
and UT1, y, m, and d are, respectively, the Universal Time, year, month, and day of
interest. JDE results from
DT
JDE ¼ JD þ
86400
and DT may be computed (in seconds) by means of the polynomial approximations

due to Espenak and Meeus [37]
DT ¼ 63:86 þ 0:3345s 0:060374 s2 þ 0:0017275 s3 þ 0:000651814 s4

þ 0:00002373599 s5 ðwith year between 1986 and 2005Þ
DT ¼ 62:92 þ 0:32217 s þ 0:005589 s2 ðwith year between 2005 and 2050Þ
DT ¼ 20 þ 32½ðy 1820Þ=1002 0:5629ð2150 yÞ
ðwith year between 2050 and 2150Þ
where
month 0:5
y ¼ year þ
12
s ¼ y 2000
Alternatively, DT may be obtained by means of the IERS tables [38].

In accordance with Kaplan [31], the following three Eulerian angles (shown in
the following figure, due to the courtesy of the National Geospatial-Intelligence
Agency [39], are needed to compute the precession matrix:
f ¼ 200 :650545 þ 230600 :083227 T þ 000 :2988499 T 2 þ 000 :01801828 T 3

000 :000005971 T 4 000 :0000003173 T 5
z ¼ 200 :650545 þ 230600 :077181 T þ 100 :0927348 T 2 þ 000 :01826837 T 3

000 :000028596 T 4 000 :0000002904 T 5
h ¼ 200400 :191903 T 000 :4294934 T 2 000 :04182264 T 3 000 :000007089 T 4

000 :0000001274 T 5
3.4 The Change of Co-ordinates Due to Precession 363
The precession matrix P(t), computed by using the three Eulerian angles f, z, and
h indicated above, is compliant with the IAU 2000A model [31].
Generally speaking, a rotation matrix (P for precession, N for nutation, etc.) is
the result of the application of the following three elementary rotation matrices
2 3 2 3
1 0 0 cos b 0 sin b
6 7 6 7
R1 4 0 cos a sin a 5 R2 4 0 1 0 5
0 sin a cos a sin b 0 cos b
2 3
cos c sin c 0
6 7
R3 4 sin c cos c 05
0 0 1
where the subscripts indicate the axis of rotation (1 for a rotation about x, 2 for a
rotation about y, and 3 for a rotation about z), as shown in the following figure.
By using the three matrices indicated above, each transformation can be written
as follows
2 3 2 3
xi x
4 yi 5 ¼ R i 4 y 5
zi z
In accordance with Kaplan [31], the complete precession matrix P(t) results from
the following sequence of rotations
PðtÞ ¼ R3 ðzÞR2 ðhÞR3 ðfÞ
By executing this matrix product, the precession matrix P(t) results

cos z cos h cos f sin z sin f cos z cos h sin f sin z cos f cos z sin h

sin z cos h cos f þ cos z sin f sin z cos h sin f þ cos z cos f sin z sin h

sin h cos f sin h sin f cos h
The matrix P(t) may also be computed by means of the ecliptic precession angles
w (luni-solar precession), x (inclination of moving equator on fixed ecliptic), v
(planetary precession), and e0J2000 (obliquity of ecliptic at J2000.0):
w ¼ 503800 :481507 T 100 :0790069 T 2 000 :00114045 T 3 þ 000 :000132851 T 4

000 :0000000951 T 5
x ¼ e0J2000 000 :025754 T þ 000 :0512623 T 2 000 :00772503 T 3 000 :000000467 T 4
þ 000 :0000003337 T 5
v ¼ 1000 :556403T 200 :3814292 T 2 000 :00121197 T 3 þ 000 :000170663 T 4
000 :0000000560 T 5
3.4 The Change of Co-ordinates Due to Precession 365
[31], where T has the same meaning as that shown above (i.e. the time, expressed in
Julian centuries, elapsed between the Julian ephemeris day of interest and J2000.0),
and e0J2000 = 23°26′21″.406 = 84381″.406 is the obliquity of the ecliptic at J2000.0
with respect to the mean equator at J2000.0. By so doing, the complete precession
matrix P(t), which is also compliant with the IAU-2000A-model, results from the
following sequence of rotations [31]:
PðtÞ ¼ R3 ðvÞR1 ðxÞR3 ðwÞR1 ðe0J2000 Þ
By executing the matrix product indicated above, the nine entries pij of the 3 3
matrix P(t) {pij} result
p11 ¼ cos v cos w þ sin w sin v cos x

p12 ¼ cos v sin w cos e0J2000 þ sin v cos x cos w cos e0J2000 þ sin e0J2000 sin v sin x
p13 ¼ cos v sin w sin e0J2000 þ sin v cos x cos w sin e0J2000 cos e0J2000 sin v sin x
p21 ¼ sin v cos w þ sin w cos v cos x
p22 ¼ sin v sin w cos e0J2000 þ cos v cos x cos w cos e0J2000 þ sin e0J2000 cos v sin x
p23 ¼ sin v sin w sin e0J2000 þ cos v cos x cos w sin e0J2000 cos e0J2000 cos v sin x
p31 ¼ sin w sin x
p32 ¼ sin x cos w cos e0J2000 cos x sin e0J2000
p33 ¼ sin x cos w sin e0J2000 þ cos x cos e0J2000
The method described above is called the “canonical 4-rotation method”.

The remaining rotation matrices (due to, respectively, nutation, sidereal time,
and polar motion) will be discussed in the following sections.
3.5 The Change of Co-ordinates Due to Nutation
The purpose of the present section is to show how to determine the rotation matrix
N(t) which takes account of nutation and the consequent changes of the ecliptic
with time. When the nutation matrix N(t) is known, it is possible to convert a
position vector in the mean-equinox-of-date system to the corresponding vector in
the true-equinox-of-date system.
Let rECImod be a position vector whose components are known in the Earth
Centred Inertial mean-equinox-of-date system. Let rECItod be the same vector whose
components are to be determined in the Earth Centred Inertial true-equinox-of-date
system. The components of rECItod result from those of rECImod by means of the
rECItod ¼ NðtÞrECImod
Let e0J2000 = 23°26′21″.406 [31] be the obliquity of the ecliptic at J2000.0 with
respect to the mean equator at J2000.0. The mean obliquity of the ecliptic of date
(e0D) with respect to the mean equator of date is
e0D ¼ e0J2000 4600 :836769 T 000 :0001831 T 2 þ 000 :00200340 T 3

000 :000000576 T 4 000 :0000000434 T 5
[31], where T = (JDE−2451545.0)/36525 is the time (expressed in Julian centuries)

elapsed between the Julian ephemeris day (JDE) of interest and J2000.0. As has
been shown in the preceding section and also in Chap. 2, Sect. 2.4, JDE is the JD
corresponding to an instant of time measured in the Terrestrial Timescale.
Meeus [40] suggests to use the following expression, due to Laskar [41], where
t is the time measured in units of 10000 Julian years from J2000.0 (i.e t = T/100):
e0D ¼ e0J2000 468000 :93 t 100 :55 t2 þ 199900 :25 t3 5100 :38 t4 24900 :67 t5
3900 :05 t6 þ 700 :12 t7 þ 2700 :87 t8 þ 500 :79 t9 þ 200 :45 t10
The accuracy of the expression given above is estimated at 0″.01 after

1000 years (i.e. between AD 1000 and 3000), and a few seconds of arc after
10000 years [40].
Meeus also points out that the same expression is valid only over a period of
10000 years on each side of J2000.0, that is, for |t| < 1.
Nutation has two components: one (Dw) is along the ecliptic and is called
nutation in longitude; the other (De) is perpendicular to the ecliptic and is called
nutation in obliquity. The five fundamental arguments used to compute Dw and De
are the Delaunay variables. Following Seidelmann [42], they are:
Mean longitude of the Moon minus the mean longitude of the Moon perigee:
l ¼ 134 570 4600 :733 þ 477198 520 0200 :633 T þ 3100 :310 T 2 þ 000 :064 T 3
Mean longitude of the Sun minus the mean longitude of the Sun perigee:
l0 ¼ 357 310 3900 :804 þ 35999 030 0100 :224 T 000 :577 T 2 000 :012 T 3
Mean longitude of the Moon minus the mean longitude of the Moon node:
F ¼ 93 160 1800 :877 þ 483202 010 0300 :137 T 1300 :257T 2 þ 000 :011 T 3
Mean elongation of the Moon from the Sun:
D ¼ 297 510 0100 :307 þ 445267 060 4100 :328 T 600 :891 T 2 þ 000 :019T 3
Longitude of the mean ascending node of the Moon orbit on the ecliptic,
measured from the mean equinox of date:
3.5 The Change of Co-ordinates Due to Nutation 367
X ¼ 125 020 4000 :280 1934 080 1000 :539 T þ 700 :455 T 2 þ 000 :008 T 3
When the five fundamental arguments indicated above have been computed as a
function of time, the nutations in longitude (Dw) and in obliquity (De) are obtained
by making the sum of the terms contained in the table “Nutation in longitude and
obliquity referred to mean ecliptic of date”, which will be given below, where the
coefficients are given in units of 0″.0001.
These terms are those of the 1980 International Astronomical Union
(IAU) theory of nutation [42]. The argument a of each sine (for Dw) and cosine (for
De) function is a linear combination of the five fundamental arguments l, l′, F, D,
and X. For example, the argument of row No. 2 of the table shown below is
0l þ 0l0 þ 0F þ 0D þ 2X ¼ 2X
When no great accuracy is required, only the periodic terms with the largest
coefficients can be used.
According to Seidelmann [43], it is possible to compute the nutations in lon-
gitude and obliquity to about 1″ by using the following expressions
Dw ¼ 0 :0048 sinð125 :0 0 :05295 dÞ 0 :0004 sinð200 :9 þ 1 :97129 dÞ

De ¼ þ 0 :0026 sinð125 :0 0 :05295 dÞ þ 0 :0002 sinð200 :9 þ 1 :97129 dÞ
where d is the number of days from 2451445.0.

Now, remembering that the mean obliquity of the ecliptic of date (e0D) with
respect to the mean equator of date is (in seconds of arc)
e0D ¼ e0J2000 4600 :836769T 000 :0001831T 2 þ 000 :00200340T 3

000 :000000576T 4 000 :0000000434T 5
the true obliquity of the ecliptic of date (eD) with respect to the equator of date is
eD ¼ e0D þ De
In accordance with Kaplan [31], the nutation matrix N(t) results from the fol-
lowing sequence of rotations
NðtÞ ¼ R1 ðeD ÞR3 ðDwÞR1 ðe0D Þ
By executing the matrix product indicated above, the nine entries nij of the 3 3
nutation matrix N {nij} result
n11 ¼ cos Dw
n12 ¼ sin Dw cos e0D
n13 ¼ sin Dw sin e0D
n21 ¼ cos eD sin Dw
n22 ¼ cos eD cos Dw cos e0D þ sin eD sin e0D
n23 ¼ cos eD cos Dw sin e0D sin eD cos e0D
n31 ¼ sin eD sin Dw
n32 ¼ sin eD cos Dw cos e0D cos eD sin e0D
n33 ¼ sin eD cos Dw sin e0D þ cos eD cos e0D
Nutation in longitude and obliquity referred to mean ecliptic of date [42].
Row l l′ F D X Dw (coefficient of De (coefficient of

sin ai) cos ai)
1 0 0 0 0 1 −171996 − 174.2 T 92025 + 8.9 T
2 0 0 0 0 2 2062 + 0.2 T −895 + 0.5 T
3 −2 0 2 0 1 46 + 0.0 T −24 + 0.0 T
4 2 0 −2 0 0 11 + 0.0 T 0 + 0.0 T
5 −2 0 2 0 2 −3 + 0.0 T 1 + 0.0 T
6 1 −1 0 −1 0 −3 + 0.0 T 0 + 0.0 T
7 0 −2 2 −2 1 −2 + 0.0 T 1 + 0.0 T
8 2 0 −2 0 1 1 + 0.0 T 0 + 0.0 T
9 0 0 2 −2 2 −13187 − 1.6 T 5736 − 3.1 T
10 0 1 0 0 0 1426 − 3.4 T 54 − 0.1 T
11 0 1 2 −2 2 −517 + 1.2 T 224 − 0.6 T
12 0 −1 2 −2 2 217 − 0.5 T −95 + 0.3 T
13 0 0 2 −2 1 129 + 0.1 T −70 + 0.0 T
14 2 0 0 −2 0 48 + 0.0 T 1 + 0.0 T
15 0 0 2 −2 0 −22 + 0.0 T 0 + 0.0 T
16 0 2 0 0 0 17 − 0.1 T 0 + 0.0 T
17 0 1 0 0 1 −15 + 0.0 T 9 + 0.0 T
18 0 2 2 −2 2 −16 + 0.1 T 7 + 0.0 T
19 0 −1 0 0 1 −12 + 0.0 T 6 + 0.0 T
20 −2 0 0 2 1 −6 + 0.0 T 3 + 0.0 T
21 0 −1 2 −2 1 −5 + 0.0 T 3 + 0.0 T
22 2 0 0 −2 1 4 + 0.0 T −2 + 0.0 T
23 0 1 2 −2 1 4 + 0.0 T −2 + 0.0 T
24 1 0 0 −1 0 −4 + 0.0 T 0 + 0.0 T
25 2 1 0 −2 0 1 + 0.0 T 0 + 0.0 T
26 0 0 −2 2 1 1 + 0.0 T 0 + 0.0 T
(continued)
(continued)
sin ai) cos ai)
27 0 1 −2 2 0 −1 + 0.0 T 0 + 0.0 T
28 0 1 0 0 2 1 + 0.0 T 0 + 0.0 T
29 −1 0 0 1 1 1 + 0.0 T 0 + 0.0 T
30 0 1 2 −2 0 −1 + 0.0 T 0 + 0.0 T
31 0 0 2 0 2 −2274 − 0.2 T 977 − 0.5 T
32 1 0 0 0 0 712 + 0.1 T −7 + 0.0 T
33 0 0 2 0 1 −386 − 0.4 T 200 + 0.0 T
34 1 0 2 0 2 −301 + 0.0 T 129 − 0.1 T
35 1 0 0 −2 0 −158 + 0.0 T −1 + 0.0 T
36 −1 0 2 0 2 −158 + 0.0 T −1 + 0.0 T
37 0 0 0 2 0 63 + 0.0 T −2 + 0.0 T
38 1 0 0 0 1 63 + 0.1 T −33 + 0.0 T
39 −1 0 0 0 1 −58 − 0.1 T 32 + 0.0 T
40 −1 0 2 2 2 −59 + 0.0 T 26 + 0.0 T
41 1 0 2 0 1 −51 + 0.0 T 27 + 0.0 T
42 0 0 2 2 2 −38 + 0.0 T 16 + 0.0 T
43 2 0 0 0 0 29 + 0.0 T −1 + 0.0 T
44 1 0 2 −2 2 29 + 0.0 T −12 + 0.0 T
45 2 0 2 0 2 −31 + 0.0 T 13 + 0.0 T
46 0 0 2 0 0 26 + 0.0 T −1 + 0.0 T
47 −1 0 2 0 1 21 + 0.0 T −10 + 0.0 T
48 −1 0 0 2 1 16 + 0.0 T −8 + 0.0 T
49 1 0 0 −2 1 −13 + 0.0 T 7 + 0.0 T
50 −1 0 2 2 1 −10 + 0.0 T 5 + 0.0 T
51 1 1 0 −2 0 −7 + 0.0 T 0 + 0.0 T
52 0 1 2 0 2 7 + 0.0 T −3 + 0.0 T
53 0 −1 2 0 2 −7 + 0.0 T 3 + 0.0 T
54 1 0 2 2 2 −8 + 0.0 T 3 + 0.0 T
55 1 0 0 2 0 6 + 0.0 T 0 + 0.0 T
56 2 0 2 −2 2 6 + 0.0 T −3 + 0.0 T
57 0 0 0 2 1 −6 + 0.0 T 3 + 0.0 T
58 0 0 2 2 1 −7 + 0.0 T 3 + 0.0 T
59 1 0 2 −2 1 6 + 0.0 T −3 + 0.0 T
60 0 0 0 −2 1 −5 + 0.0 T 3 + 0.0 T
61 1 −1 0 0 0 5 + 0.0 T 0 + 0.0 T
62 2 0 2 0 1 −5 + 0.0 T 3 + 0.0 T
63 0 1 0 −2 0 −4 + 0.0 T 0 + 0.0 T
64 1 0 −2 0 0 4 + 0.0 T 0 + 0.0 T
65 0 0 0 1 0 −4 + 0.0 T 0 + 0.0 T
(continued)
(continued)
sin ai) cos ai)
66 1 1 0 0 0 −3 + 0.0 T 0 + 0.0 T
67 1 0 2 0 0 3 + 0.0 T 0 + 0.0 T
68 1 −1 2 0 2 −3 + 0.0 T 1 + 0.0 T
89 −1 −1 2 2 2 −3 + 0.0 T 1 + 0.0 T
70 −2 0 0 0 1 −2 + 0.0 T 1 + 0.0 T
71 3 0 2 0 2 −3 + 0.0 T 1 + 0.0 T
72 0 −1 2 2 2 −3 + 0.0 T 1 + 0.0 T
73 1 1 2 0 2 2 + 0.0 T −1 + 0.0 T
74 −1 0 2 −2 1 −2 + 0.0 T 1 + 0.0 T
75 2 0 0 0 1 2 + 0.0 T −1 + 0.0 T
76 1 0 0 0 2 −2 + 0.0 T 1 + 0.0 T
77 3 0 0 0 0 2 + 0.0 T 0 + 0.0 T
78 0 0 2 1 2 2 + 0.0 T −1 + 0.0 T
79 −1 0 0 0 2 1 + 0.0 T −1 + 0.0 T
80 1 0 0 −4 0 −1 + 0.0 T 0 + 0.0 T
81 −2 0 2 2 2 1 + 0.0 T −1 + 0.0 T
82 −1 0 2 4 2 −2 + 0.0 T 1 + 0.0 T
83 2 0 0 −4 0 −1 + 0.0 T 0 + 0.0 T
84 1 1 2 −2 2 1 + 0.0 T −1 + 0.0 T
85 1 0 2 2 1 −1 + 0.0 T 1 + 0.0 T
86 −2 0 2 4 2 −1 + 0.0 T 1 + 0.0 T
87 −1 0 4 0 2 1 + 0.0 T 0 + 0.0 T
88 1 −1 0 −2 0 1 + 0.0 T 0 + 0.0 T
89 2 0 2 −2 1 1 + 0.0 T −−1 + 0.0 T
90 2 0 2 2 2 −1 + 0.0 T 0 + 0.0 T
91 1 0 0 2 1 −1 + 0.0 T 0 + 0.0 T
92 0 0 4 −2 2 1 + 0.0 T 0 + 0.0 T
93 3 0 2 −2 2 1 + 0.0 T 0 + 0.0 T
94 1 0 2 −2 0 −1 + 0.0 T 0 + 0.0 T
95 0 1 2 0 1 1 + 0.0 T 0 + 0.0 T
96 −1 −1 0 2 1 1 + 0.0 T 0 + 0.0 T
97 0 0 −2 0 1 −1 + 0.0 T 0 + 0.0 T
98 0 0 2 −1 2 −1 + 0.0 T 0 + 0.0 T
99 0 1 0 2 0 −1 + 0.0 T 0 + 0.0 T
100 1 0 −2 −2 0 −1 + 0.0 T 0 + 0.0 T
101 0 −1 2 0 1 −1 + 0.0 T 0 + 0.0 T
102 1 1 0 −2 1 −1 + 0.0 T 0 + 0.0 T
103 1 0 −2 2 0 −1 + 0.0 T 0 + 0.0 T
104 2 0 0 2 0 1 + 0.0 T 0 + 0.0 T
(continued)
(continued)
sin ai) cos ai)
105 0 0 2 4 2 −1 + 0.0 T 0 + 0.0 T
106 0 1 0 1 0 1 + 0.0 T 0 + 0.0 T
In order to use the table shown above, the argument (ai) of the sine and cosine
functions is computed, for each row i, as follows
ai ¼ ai l þ bi l 0 þ c i F þ di D þ e i X
where the fundamental argument multipliers ai, bi, ci, di, and ei are given in the
second, third, fourth, fifth, and sixth columns of the table. Then, the nutations in
longitude (Dw) and in obliquity (De) are computed as follows
X
n X
n
Dw ¼ si sin ai De ¼ ci cos ai
i¼1 i¼1
where n = 106 is the number of rows of the table, si and ci are the coefficients of,
respectively, the sine function (for Dw) and the cosine function (for De), which are
given in, respectively, the seventh and eighth columns of the table.
According to McCarthy [44], observations performed by means of Very Long
Baseline Interferometry (VLBI) and Lunar Laser Ranging (LLR) techniques have
shown deficiencies in the IAU Precession and in the IAU 1980 Theory of Nutation.
However, these models are kept as part of IERS standards, and the observed dif-
ferences (dDw and dDe, equivalent to dPsi and dEps in the IERS bulletins, see for
example [45]) with respect to the conventional celestial pole position defined by the
models are monitored and reported by the IERS as “celestial pole offsets”. Using
these offsets, the corrected nutations in, respectively, longitude (Dw*) and obliquity
(De*) are given by
Dw ¼ Dw þ dDw
De ¼ De þ dDe
According to McCarthy [46], the observed differences dDw and dDe can be
taken into account by replacing the nutation matrix N shown above (computed
according to the 1980 nutation theory) with another matrix N* such that

N ¼ RN
where
2 3
1 dDw cos D dDw sin D
4 dDw cos D
R 1 dDx 5
dDw sin D dD 1
This makes it possible to take account of the celestial pole offsets given in the
IERS bulletins. The celestial pole offsets will be considered again in Sect. 3.9.
More recently, the Circular No. 179 of the IAU [31] gives a table, called IAU
2000A nutation series, containing as many as 14 fundamental argument multipliers
(instead of the 5 multipliers of the 1980 theory) and 1365 rows (instead of the 106
rows of the 1980 theory). The new nutation model described in the IAU Circular
No. 179 is shown below.
According to the IAU 2000A nutation model, the nutations in longitude (Dw)
and in obliquity (De) are obtained by evaluating the following trigonometric series
N h
X i
Dw ¼ Si þ Sbi sin Ui þ Ci0 cos Ui
i¼1
N h
X i
D ¼ Ci þ c
Ci sin Ui þ S0i cos Ui
i¼1
where in each term

X
K
Ui ¼ Mij uj ðT Þ
j¼1
As has been shown above, N = 1365, K = 14, and T = (JDE−2451545.0)/36525

is the time (expressed in Julian centuries) elapsed between the Julian ephemeris day
(JDE) of interest and J2000.0. Strictly speaking, T denotes, in the preceding
expression, the number of Julian centuries of TDB (which stands for Barycentric
Dynamical Time) since the 1st of January 2000, 12 h TDB. However, Kaplan [31]
points out that the error, committed in using TDB instead of Terrestrial Time (TT),
is negligible. By the way, TDB is the same as TT, except for relativistic corrections
to move the origin to the Solar System barycentre. These corrections, of about
1.6 ms, are periodic with an average value of zero. The dominant terms in these
corrections have annual and semi-annual periods. Fisher [47] gives an approximate
relation (in seconds) between TDB and TT:
TDB ¼ TT þ 0:001658 sinðgÞ þ 0:000014 sinð2gÞ
where g = 357.53 + 0.9856003 (JD−2451545.0) is the mean anomaly (in degrees)

of the Earth. Further information on the matter can be found in Ref. [31].
The first eight of the fourteen fundamental arguments u1(T), u2(T), …, u14(T) are
the mean heliocentric ecliptic longitudes of the planets from Mercury to Neptune:
u1 ¼ 90810300 :259872 þ 53810162800 :688982T

u2 ¼ 65512700 :283060 þ 21066413600 :433548T
u3 ¼ 36167900 :244588 þ 12959774200 :28342900 T
u4 ¼ 127955800 :798488 þ 6890507700 :493988T
u5 ¼ 12366500 :467464 þ 1092566000 :377991T
u6 ¼ 18027800 :799480 þ 439960900 :855732T
u7 ¼ 113059800 :018396 þ 154248100 :193933T
u8 ¼ 109565500 :195728 þ 78655000 :320744T
Kaplan [31] points out that in certain applications, it may be necessary to convert
the angles given above to radians in the range [0, 2p]. The ninth, u9(T), of the
fundamental arguments is an approximation to the general precession in longitude:
u9 ¼ 5028:8200T þ 1:112022T 2
The last five arguments, from u10(T) to u14(T), are the same luni-solar angles as
those used in the 1980 IAU theory of nutation [42], namely l, l′, F, D, and X, the
only difference being the way of expressing these angles (given below in seconds of
arc) as a function of time:
u10 ¼ l ¼ 485868:249036 þ 1717915923:2178T þ 31:8792T 2 þ 0:051635T 3

0:00024470T 4
u11 ¼ l0 ¼ 1287104:79305 þ 129596581:0481T 0:5532T 2 þ 0:000136T 3
0:00001149T 4
u12 ¼ F ¼ 335779:526232 þ 1739527262:8478T 12:7512T 2 0:001037T 3
þ 0:00000417T 4
u13 ¼ D ¼ 1072260:70369 þ 1602961601:2090T 6:3706T 2 þ 0:006593T 3
0:00003169T 4
u14 ¼ X ¼ 450160:398036 6962890:5431T þ 7:4722T 2 þ 0:007702T 3
0:00005939T 4
In order to evaluate the IAU 2000A nutation series for a given date t of interest,
it is necessary to compute first the fourteen fundamental arguments u1(T), u2(T),
…, u14(T). This is done only once. Then, the nutation terms are evaluated one by
one, as follows. For each term i = 1, 2, …, 1365, Ui is computed by means of
X
14
Ui ¼ Mij uj ðT Þ
j¼1
where the fourteen fundamental argument multipliers Mij (j = 1, 2, …, 14) are given
in pages 88–103 of [31]. Then, the cosine and sine components for each term are
evaluated by means of the expressions
1365 h
X i
Dw ¼ Si þ Sbi sin Ui þ Ci0 cos Ui
i¼1
1365 h
X i
D ¼ Ci þ c
Ci sin Ui þ S0i cos Ui
i¼1
The table named above contains 1365 rows and 14 columns and occupies 16
pages. Therefore, it is not given here. The reader can find it in Ref. [31].
In practice, the IAU 2000A nutation series is evaluated by means of an ad hoc
piece of software (SOFA tools for Earth attitude), which is described in Ref. [48].
3.6 The Change of Co-ordinates Due to the Rotation

of the Earth
The present section shows how to determine the spin matrix S(t), also called
sidereal-time matrix, which takes account of the rotation of the Earth around its
axis. This is the transformation which induces the largest numerical change in the
position vector components to which it is applied. The application of the transpose
matrix [S(t)]T transforms such components in the non-rotating true-of-date system
to those in the Earth-fixed system. A true-of-date reference system is defined by
means of the true equator and equinox on the date of interest. After the transfor-
mation, the X-axis of the true-of-date non-rotating system gets aligned with the x-
axis of the system co-rotating with the Earth. Let rECItod be a position vector whose
components are known in the Earth Centred Inertial true-equinox-of-date system.
Let rECEFw/oPM be the same vector whose components are to be determined in the
Earth Centred Earth Fixed system not corrected for polar motion (also called
Pseudo Earth-Fixed, PEF). The components of rECEFw/oPM result from those of
rECItod by means of the following expression
rECEFw=oPM ¼ ½SðtÞT rECItod
As has been anticipated in Sect. 3.2, this calculation requires the knowledge of
the Greenwich apparent sidereal time (hG), which is the Greenwich mean sidereal
time (hG0) corrected for the shift in the position of the vernal equinox due to
nutation. According to McCarthy [46], hG results from
hG ¼ hG0 þ Dw cos e0D þ 000 :00264 sin X þ 000 :000063 sinð2XÞ

3.6 The Change of Co-ordinates Due to the Rotation of the Earth 375
where Dw is the nutation in longitude, which can be computed as shown in

Sect. 3.5, e0D is the mean obliquity of the ecliptic of date, and X is the mean
longitude of the ascending node of the lunar orbit (see Sect. 3.5). The corrective
terms after hG0, which are due to nutations, are called the equation of the equinoxes
or the nutation in right ascension. The Greenwich mean sidereal time (hG0), in turn,
measures the Earth rotation with respect to the distant stars and takes account of the
smoothly varying part of the change of direction of Earth rotation axis due to
precession. The Greenwich mean sidereal time is the hour angle of the average
position of the vernal equinox, where the short-term motions of the equinox, due to
nutation, are not taken into account. According to McCarthy [46], in order to
compute hG0, we first compute hG0 at 0 h UT1, by using the following expression
due to Aoki et al. [49, p. 360, Eq. 13]:
hG0 ðof 0h UT1Þ ¼ 24110s :54841 þ 8640184s :812866Tu0 þ 0s :093104Tu0 2

6s :2 106 Tu0 3
where T′u = d′u/36525 and d′u is the number of days elapsed since the 1st of January
2000, 12 h UT1, taking on values ±0.5, ±1.5, …, and then compute hG0, as
follows
hG0 ¼ hG0 ðof 0h UT1Þ þ r½ðUT1 UTC) þ UTC
where r, which is the ratio of universal to sidereal time, as given by Aoki et al. [49,
p. 361, Eq. 19]:
r ¼ 1:002737909350795 þ 5:9006 1011 Tu0 5:9 1015 Tu0 2
and the value of UT1 − UTC is the IERS value.

After computing hG0 as shown above, hG results from the preceding expression
hG ¼ hG0 þ Dw cos e0D þ 000 :00264 sin X þ 000 :000063 sinð2XÞ
This method reflects the IERS standards of 1992 for computing the Greenwich
apparent sidereal time (hG).
More recently, the method for computing hG has been modified as will be shown
below. Following Kaplan [31], the Earth rotation angle (h) is computed as follows
h ¼ 0:7790572732640 þ 0:00273781191135448DU þ frac[JD(UT1Þ
where DU is the number of UT1 days from 1 January 2000, 12 h UT1 (in other
words, DU = JD(UT1)−2451545.0), and frac[JD(UT1)] is the fractional part of the
UT1 Julian date (i.e. JD(UT1) modulus 1.0). When the value of h is known, the
Greenwich mean sidereal time in seconds results from
1
hG0 ¼ 86400h þ ð0:014506 þ 4612:156534T þ 1:3915817T 2
15
0:00000044T 3 0:000029956T 4 0:0000000368T 5 Þ
where T is the number of centuries of TDB (or equivalently, for this purpose, of TT)
from J2000.0 (in other words, T = JD(TDB or TT)−2451545.0)/36525).
The polynomial in parentheses is the accumulated precession of the equinox in
right ascension, expressed in arcseconds. Kaplan [31] points out that two timescales
are required to compute hG0: in the “fast term”, h is a function of UT1, whereas, in
the remaining terms, T is expressed in TDB or TT.
When the value of hG0 is known, the Greenwich apparent mean sidereal time
(hG) is obtained in seconds by adding the equation of the equinoxes to hG0, as
follows
!
hG ¼ hG0 þ
15
As has been shown above, the addition of the second term on the right-hand side
of the preceding expression is due to the motion of the equinox because of nutation.
The new theory of nutation, shown at the end of sect. 3.5, takes account of the
so-called complementary terms and expresses the equation of the equinoxes, in
arcseconds, as follows
e! ¼ Dw cos e0D þ 0:00264096 sin X þ 0:00006352 sinð2XÞ

þ 0:00001175 sinð2F 2D þ 3XÞ þ 0:00001121 sinð2F 2D þ XÞ
0:00000455 sinð2F 2D þ 2XÞ þ 0:00000202 sinð2F þ 3XÞ
þ 0:00000198 sinð2F þ XÞ 0:00000172 sinð3XÞ 0:000000 sin X þ
where, again, Dw is the nutation in longitude (in arcseconds), e0D is the mean
obliquity of the ecliptic of date, and F, D, and X are the fundamental luni-solar
arguments. Finally, in accordance with Kaplan [31], McCarthy [46], and Vallado
et al. [50], the sidereal-rotation matrix S(t) results from
SðtÞ ¼ R3 ðhG Þ
Hence, remembering (see Sect. 3.4) the definition of R3, there results
2 3 2 3
cos c sin c 0 cos hG sin hG 0
R3 ðcÞ 4 sin c cos c 05 SðtÞ 4 sin hG cos hG 05
0 0 1 0 0 1
3.6 The Change of Co-ordinates Due to the Rotation of the Earth 377
Hence, the time derivative, S′(t), and the time derivative transpose, [S′(t)]T, of
the sidereal-rotation matrix, S(t), are respectively
2 3
xE sin hG xE cos hG 0
6 7
S0 ðtÞ 4 xE cos hG xE sin hG 05
0 0 0
2 3
0 T 6 7
½S ðtÞ 4 xE cos hG xE sin hG 05
0 0 0
where xE is the true angular rotation of the Earth, whose magnitude xE, according
to the IERS [51], is given by

LOD
xE ¼ xN 1
t
where xN = 72921151.467064 10−12 radians/second is the nominal rotation rate

of the Earth (corresponding to the mean rotation rate of the epoch 1820),
t = 86400 s TAI, and LOD is the excess of length of day, that is, the difference
between the astronomically determined duration of the day and 86400 s. Tables
giving the value (expressed in milliseconds) of LOD for the day of interest can be
found in Ref. [52].
3.7 The Change of Co-ordinates Due to Polar Motion
The present section shows how to determine the wobble rotation matrix W(t), also
called polar-motion matrix, which accounts for the variation of the exact location of
the north pole of the Earth. The application of the transpose matrix [W(t)]T
transforms a position vector whose components are known in the Earth Centred
Earth Fixed system not corrected for polar motion to the same vector in the Earth
Centred Earth Fixed system corrected for polar motion, by using the actual location
of the IERS Reference Pole.
Let rECEFw/oPM be a position vector whose components are known in the Earth
Centred Earth Fixed system not corrected for polar motion. Let rECEFwPM be the
same vector whose components are to be determined in the Earth Centred Earth
Fixed system corrected for polar motion.
The components of rECEFwPM result from those of rECEFw/oPM as follows
rECEFwPM ¼ ½WðtÞT rECEFw=oPM

The motion of the rotation axis of the Earth with respect to the crust is repre-
sented by two angles, xP and yP, which are shown in the following figure.
They are the angles between the Celestial Ephemeris Pole (CEP) of date and the
IERS Reference Pole (IRP). The latter, in turn, is the mean location of the pole as
defined by agreements reached by international committees on the basis of actual
observations.
In accordance with McCarthy [53], the co-ordinate system used for polar motion
is a xy-plane, whose x-axis is in the direction of the IERS Reference Meridian
(which, in turn, is 5.31 arcseconds east of Airy’s transit circle or 102.5 m at the
latitude of the Royal Observatory, Greenwich) and points to the south, and whose y-
axis is in the direction 90° west longitude. Thus, xP and yP are the angular
co-ordinates of the Celestial Ephemeris Pole relative to the IERS Reference Pole.
The CEP differs from the instantaneous rotation axis by quasi-diurnal terms with
amplitudes under 0.01 arcseconds.
In accordance with McCarthy [46] and Vallado et al. [50], the polar-motion
matrix W(t) results from two rotations, as follows
WðtÞ ¼ R1 ðyP ÞR2 ðxP Þ
Since the matrices R1 and R2, defined in Sect. 3.4, are

2 3 2 3
1 0 0 cos b 0 sin b
R1 ðaÞ 4 0 cos a sin a 5 R2 ðbÞ 4 0 1 0 5
0 sin a cos a sin b 0 cos b
3.7 The Change of Co-ordinates Due to Polar Motion 379
then there results

2 32 3
1 0 0 cos xP 0 sin xP
6 76 7
R1 ðyP ÞR2 ðxP Þ ¼ 4 0 cos yP sin yP 54 0 1 0 5
0 sin yP cos yP sin xP 0 cos xP
2 3 2 3
cos xP 0 sin xP 1 0 xP
6 7 6 7
¼ 4 sin yP sin xP cos yP sin yP cos xP 5 4 0 1 yP 5
cos yP sin xP sin yP cos yP cos xP xP yP 1
This approximation is justified because xP and yP are very small angles (the
maximum amplitude observed hitherto is about 0″.3). Then, in accordance with
McCarthy [53], the resulting polar-motion matrix is
2 3
1 0 xP
Wð t Þ 4 0 1 yP 5
xP yP 1
where xP and yP are the angular co-ordinates, expressed in radians, of the Celestial
Ephemeris Pole relative to the IERS Reference Pole.
3.8 The Fundamental Reference Systems
Summarising what has been shown in the preceding sections, the procedure for
converting a position vector r in the Earth-Centred Inertial J2000 system (denoted
with the subscript J2000mod and based on the mean equinox of J2000) to the
corresponding vector in the Earth-Centred Earth-Fixed system (denoted with the
subscript ECEFwPM and corrected for polar motion) is the following
rJ2000mod ¼ PðtÞNðtÞSðtÞWðtÞrECEFwPM
where t is the observation epoch in Terrestrial Time (TT), P(t) is the precession
matrix at epoch t, N(t) is the nutation matrix at epoch t, S(t) is the spin matrix at
epoch t due to the daily rotation of the Earth around the celestial pole, and W(t) is
the wobble (or polar-motion) matrix describing the position of the celestial pole in
the Earth-Centred Inertial J2000 system at epoch t.
The integration of the differential equations giving the position and velocity
vectors of an Earth-orbiting satellite as a function of the forces acting on it requires
the preliminary definition of an inertial reference system. This is because the
equations of motion, if they were written in a non-inertial reference system, would
require the addition of fictitious forces. However, several of the perturbing forces
acting on a satellite (the gravitational force for one) have been expressed so far in a
reference system co-rotating with the Earth in its diurnal motion. Other perturbing
forces to be considered in the following sections will also be expressed, for con-
venience, in the same co-rotating reference system. In addition, the observations
performed in stations located on the surface of the Earth also refer to the same
co-rotating reference system. Hence, before integrating the equations of motion, it
is necessary to define exactly the inertial reference system, the terrestrial (i.e.
Earth-fixed) reference system, and the equations needed to transform the compo-
nents of a given vector from those computed in one of the two reference systems to
those relating to the other. The reference systems used in practice in astronomy and
the relations existing between one of them to another are established and approved
by international organisations, such as the International Astronomical Union
(IAU) and the International Union of Geodesy and Geophysics (IUGG). The object
of the present section is to describe briefly such reference systems.
The International Celestial Reference System (ICRS) is defined kinematically,
making the directions of its axes fixed with respect to distant objects of the uni-
verse. In accordance with Kaplan [31], the ICRS is a co-ordinate system whose
origin is at the Solar System barycentre and whose axis directions are effectively
defined by the adopted co-ordinates of 212 extragalactic radio sources observed by
Very Long Baseline Interferometry (VLBI). These radio sources (quasars and active
galactic nuclei) are assumed to have no observable intrinsic angular motions.
Therefore, the ICRS is a kinematically non-rotating system without an associated
epoch. However, the ICRS closely matches the conventional dynamical system
defined by the Earth mean equator and equinox of J2000.0; the alignment difference
is at the 0.02 arcsecond level, which is negligible for many applications [31]. As
will we shown in the next sections, the so-called frame-bias matrix B is required to
convert co-ordinates from the International Celestial Reference System (ICRS) to
the dynamical mean equator and equinox of J2000.0 (considered for this purpose to
be a barycentric system). The same matrix is used in the geocentric transformations
(described in the next sections) involving vectors in the GCRS (i.e. in the
Geocentric ICRS), so that they can be operated on by the conventional precession
and nutation matrices. In the barycentric case, the matrix B is used as follows
rJ2000mod ¼ BrICRS
where rJ2000mod is a position vector with respect to the dynamical mean equator and
equinox of J2000.0, and rICRS is the corresponding vector with respect to the ICRS.
In the geocentric case, rICRS is replaced by rGCRS, and rECImod is also a geocentric
vector. To this regard, Seitz and Schuh [54] point out that, if the origin of the ICRS
is shifted from the barycentre of the Solar System into the Earth centre of mass
(under consideration of relativistic effects), the system experiences slight acceler-
ations due to the motion of the Earth around the Sun. Therefore, strictly speaking,
such a Geocentric Celestial Reference System (GCRS) is no longer an inertial
system and is commonly referred to as a quasi-inertial system [54].
The International Astronomical Union (IAU) recommended that the origin of
right ascensions of the ICRS should be close to the dynamical equinox at J2000.0.
3.8 The Fundamental Reference Systems 381
The X-axis of the IERS celestial system was implicitly defined in its initial reali-
sation by adopting the mean right ascension of 23 radio sources in a group of
catalogues that were compiled by fixing the right ascension of the quasar 3C 273B
to the usual conventional value, that is, to 12h 29m 6s.6997 at J2000.0 [55].
Analyses of Lunar Laser Ranging (LLR) observations indicate that the origin of
right ascension in the ICRS is shifted from the inertial mean equinox at J2000.0 on
the ICRS reference plane by −55.4 ± 2.3 milliarcseconds (mas) (direct rotation
around the polar axis). This shift of −55.4 mas on the ICRS equator corresponds to
a shift of −14.6 mas on the mean equator of J2000.0 [55]. Kaplan [31] distinguishes
between the International Celestial Reference System (ICRS), described above, and
the International Celestial Reference Frame (ICRF), the latter being a set of
extragalactic objects whose adopted positions and uncertainties realise the ICRS
axes and give the uncertainties of the axes. ICRF is also the name of the radio
catalogue whose 212 defining sources are currently the most accurate realisation of
the ICRS [31].
According to Petit and Luzum [56], the International Terrestrial Reference
System (ITRS) is a spatial reference system co-rotating with the Earth in its diurnal
motion in space, such that the positions of points attached to the solid surface of the
Earth have co-ordinates which undergo only small variations with time, due to
geophysical effects (tectonic or tidal deformations). In particular, the ITRS fulfils
the following conditions:
(1) it is geocentric, its origin being the centre of mass for the whole Earth,
including oceans and atmosphere;
(2) the unit of length is the metre (SI);
(3) its orientation was initially given by the Bureau International de l’Heure
(BIH) orientation at 1984.0;
(4) the time evolution of the orientation is ensured by using a no-net rotation
condition with regard to horizontal tectonic motions over the whole Earth.
The ITRS was aligned close to the mean equator of 1900 and the Greenwich
meridian, for continuity with previous terrestrial reference systems. The ITRS is the
recommended system to express positions on the Earth and is realised by the
International Terrestrial Reference Frame (ITRF), which in turn is a set of physical
points with precisely determined co-ordinates in the ITRS.
3.9 The Frame-Bias Matrix
As shown in Sect. 3.8, the frame-bias matrix B is necessary to convert position

vectors from the International Celestial Reference System (ICRS) to the dynamical
mean equator and equinox of J2000.0 (considered for this purpose to be a
barycentric system). The same matrix is used in the geocentric transformations

(described in the next sections) involving vectors in the GCRS (i.e. in the
Geocentric ICRS), so that they can be operated on by the conventional precession
(P) and nutation (N) matrices, which have been described in the preceding sections.
In the geocentric case, let rGCRS be a position vector in the Geocentric ICRS
system and let rJ2000mod be the corresponding vector in the dynamical mean equator
and equinox of J2000.0 system. According to Kaplan [31], the frame-bias matrix B,
such that
rJ2000mod ¼ BrGCRS
results from the following sequence of rotations
B ¼ R1 ðdeB ÞR2 ðdwB sin e0J2000 ÞR3 ðda0 Þ
where R1(a), R2(b), and R3(c) are the three elementary matrices which have been
described in Sect. 3.4, e0J2000 = 84381″.406 is the obliquity of the ecliptic at
J2000.0 with respect to the mean equator at J2000.0, dwB = −0″.041775 and
deB = −0″.0068192 are, respectively, the longitude and obliquity of the celestial
pole, and da0 = −0″.0146 is the offset in the ICRS right ascension origin with
respect to the dynamical equinox of J2000.0, as measured in an inertial
(non-rotating) reference system. Let
n0 ¼ dwB sin e0J2000 ¼ 000 :041775 sinð8438100 :406Þ ¼ 000 :016617

g0 ¼ deB ¼ 000 :0068192
da0 ¼ 000 :0146
be the frame bias in rectangular co-ordinates. The celestial pole offsets (n0, η0, da0)
are the angular displacements between the observed position of the rotation axis of
the Earth in space and its position predicted by the models (described in the pre-
ceding sections) of precession and nutation. Using the rectangular co-ordinates n0
and η0 instead of the longitude (dwB) and obliquity (deB) of the celestial pole, the
preceding expression of the frame-bias matrix B becomes
B ¼ R1 ðg0 ÞR2 ðn0 ÞR3 ðda0 Þ
Executing this matrix product leads to the following components of B {bij}:

3.9 The Frame-Bias Matrix 383
b11 ¼ cos n0 cos da0

b12 ¼ cos n0 sin da0
b13 ¼ sin n0
b21 ¼ cos g0 sin da0 sin g0 sin n0 cos da0
b22 ¼ cos g0 cos da0 sin g0 sin n0 sin da0
b23 ¼ sin g0 cos n0
b31 ¼ sin g0 sin da0 þ cos g0 sin n0 cos da0
b32 ¼ sin g0 cos da0 þ cos g0 sin n0 sin da0
b33 ¼ cos g0 cos n0
Since the arguments of the sine and cosine functions are of the order of mag-
nitude of tenths of arcseconds, these functions can be expanded in a Maclaurin
series truncated after the first term (that is, cos x 1 and sin x x, with x in
radians). This leads to the following first-order approximation:
2 3
1 da0 n0
B 4 da0 1 g0 5
n0 g0 1
where, as shown above, the celestial pole offsets (n0, η0, da0) must be expressed in radians.
A better approximation results from setting cos x 1 − x2/2!, sin x x − x3/3!,
and neglecting all monomials whose degree is greater than 2:
1
b11 ¼ cos n0 cos da0 1 n20 =2 1 da20 =2 ¼ 1 da20 þ n20
2
b12 ¼ cos n0 sin da0 1 n20 =2 da0 da30 =6 ¼ da0

b13 ¼ sin n0 n0 n30 =6 ¼ n0

b21 ¼ cos g0 sin da0 sin g0 sin n0 cos da0 ð1 g20 =2Þ da0 da30 =6

ðg0 g30 =6Þ n0 n30 =6 1 da20 =2 ¼ da0 g0 n0

b22 ¼ cos g0 cos da0 sin g0 sin n0 sin da0 ð1 g20 =2Þ 1 da20 =2
1
ðg0 g30 =6Þ n0 n30 =6 da0 da30 =6 ¼ 1 da20 þ g20
2
b23 ¼ sin g0 cos n0 g0 g30 =6 1 n20 =2 ¼ g0

b31 ¼ sin g0 sin da0 þ cos g0 sin n0 cos da0 ðg0 g30 =6Þ da0 da30 =6

þ ð1 g20 =2Þ n0 n30 =6 1 da20 =2 ¼ g0 da0 þ n0

b32 ¼ sin g0 cos da0 þ cos g0 sin n0 sin da0 g0 g30 =6 1 da20 =2

þ 1 g20 =2 n0 n30 =6 da0 da30 =6 ¼ g0 þ n0 da0
1
b33 ¼ cos g0 cos n0 1 g20 =2 1 n20 =2 ¼ 1 n20 þ g20
2
This leads to the following second-order approximation:

2 3
1 12 da20 þ n20 da0 n0
B 4 da0 g0 n0 1 12 da20 þ g20 g0

5
n0 g0 da0 g0 þ n0 da0 1 12 n20 þ g20
3.10 The Co-ordinate Transformation, Based

on the Equinox, Between the Celestial and Terrestrial
Reference Systems
The changes of co-ordinates described in the preceding sections lead to the fol-
lowing transformation:
rGCRS ¼ BPðtÞNðtÞSðtÞWðtÞrITRS
where rGCRS is the position vector of the object of interest in the Geocentric ICRS,
rITRS is the corresponding vector in the International Terrestrial Reference System,
and B, P, N, S R3(–hG), and W R1(yP) R2(xP) are, respectively, the
frame-bias, precession, nutation, sidereal-rotation, and polar-motion matrices which
have been described in Sects. 3.4–3.9. All of these matrices, except B, depend on
time. The inverse transformation is
rITRS ¼ ½WðtÞT ½SðtÞT ½NðtÞT ½PðtÞT ½BT rGCRS
This is because B, P, N, S, and W are orthogonal matrices (the explicit

time-dependence of P, N, S, and W will henceforth be dropped for convenience).
In case of velocity vectors (v), the transformation is as follows [50]:
rECEFw=oPM ¼ ST NT PT rJ2000mod
n o
vITRS ¼ ½WT ½ST ½BPNT vGCRS xE rECEFw=oPM
¼ ½WT ½ST ½BPNT vGCRS þ ½WT ½S0 ½BPNT rGCRS

T

vGCRS ¼ ½BPN ½S ½WvITRS þ xE rECEFw=oPM
¼ ½BPN ½S ½W vITRS þ ½BPN½S0 ½WrGCRS
where [BPN] is either the single time-dependent bias-precession-nutation matrix

(which will be described in Sect. 3.11) or the matrix resulting from the product of
the tree matrices B, P, and N (which have been described in the previous sections),
depending on whether the IAU-2000A theory or the IAU-1976/1980 theory of
precession/nutation is adopted; xE is the true angular rotation of the Earth; and
3.10 The Co-ordinate Transformation, Based on the Equinox … 385
rECEFw/oPM (also denoted with rPEF) is the position vector in the Earth Centred Earth
Fixed system not corrected for polar motion.
As has been shown in Sect. 3.6, S′ and [S′]T have the following expressions:
2 3
6 7
S0 ðtÞ 4 xE cos hG xE sin hG 05
0 0 0
2 3
0 T 6 7
½S ðtÞ 4 xE cos hG xE sin hG 05
0 0 0
where hG the Greenwich apparent sidereal time.

The method described above is the classical procedure, which uses the equinox
to realise the intermediate reference system of date t. This reference system is
defined below.
Following de Viron et al. [57], let us consider the instantaneous rotation axis of
the Earth, that is, the axis of the points which are not affected by the rotation of the
Earth at any given time. This axis is not fixed in the terrestrial reference system, but
rather moves with respect to the latter as the rotation axis of the Earth moves with
respect to the crust (as a result of the polar motion described in Sects. 3.3 and 3.7).
Let us consider another axis, whose motion with respect to the terrestrial ref-
erence system is the same as that of the instantaneous rotation axis of the Earth, but
smoothed over a one-day period and consistent with the definition of polar motion
and nutation given in the preceding sections. This axis is called the Celestial
Intermediate Pole (CIP) axis, in conformity with the recent resolutions of the
International Astronomical Union (IAU) and the International Union of Geodesy
and Geophysics (IUGG). These resolutions were taken because the original defi-
nition of the Celestial Ephemeris Pole (CEP), as given in Sect. 3.7, is based upon
the IAU-1976/1980 theory of precession/nutation, which does not take account of
diurnal and higher-frequency variations in the Earth orientation. Consequently, the
XXIV General Assembly of the IAU, held in Manchester in August 2000, has
extended the concept of CEP to that one of CIP.
According to the IERS [58], the CIP is an intermediate pole separating the
motion of the pole of the ITRS in the GCRS by dividing nutation from polar motion
explicitly at the 2-day nutation period. Nutations with periods less than two days are
modelled by their equivalent polar motion.
The Celestial Intermediate Pole (CIP) is the Reference Pole of the IAU 2000A
precession/nutation theory. The CIP is the pole, whose motion is specified in the
GCRS by the motion of the Tisserand mean axis of the Earth with periods greater
than two days.
The Tisserand mean axis of the Earth, in turn, is the axis about which the total
internal angular momentum of the Earth is zero.
The motions of the CIP, described by the IAU-1976/1980 theory of precession/
nutation, are those of the Tisserand mean axis with periods greater than two days in
the GCRS.
Let x, y, and z be the Cartesian axes of the (geocentric) terrestrial reference

system (ITRS). Let X, Y, and Z be the Cartesian axes of the GCRS.
In line of principle, the transformation from the terrestrial to the celestial ref-
erence system can be done (in both the classical and the new theory) in three
successive steps:
• by means of a rotation bringing the pole of the terrestrial reference system on the
Celestial Intermediate Pole (CIP), corresponding to the motion of the CIP with
respect to the terrestrial reference system, the z-axis of the terrestrial reference
system is made to coincide with the axis corresponding to the CIP;
• by means of a rotation around the new z-axis (which is oriented in the direction
of the CIP), account is taken of the Earth rotation (for either the Greenwich
sidereal time or the Earth rotation angle, as will be shown below); and
• by means of a rotation bringing the new equatorial plane on the celestial equator,
the CIP is brought on the Z-axis of the celestial reference system.
The three steps described above cannot be taken by only one transformation, and
consequently two intermediate reference systems are needed.
One of these intermediate reference systems does rotate with the Earth, whereas
the other does not. They share, both of them, the same xy-plane, but differ one from
the other because their x-axes are separated by the Greenwich sidereal time
(GST) or the Earth rotation angle (ERA).
The three rotations described above can be chosen differently, according to the
point chosen for the intermediate pole. The classical theory (described above) and
the new theory (to be described in Sect. 3.11) define differently the two interme-
diate reference systems.
3.11 The Co-ordinate Transformation, Based

on the Non-rotating Origins, Between the Celestial
and Terrestrial Reference Systems
The transformation based on the non-rotating origins uses the same intermediate
pole (the CIP) but an intermediate reference system which differs from that which is
used in the equinox-based transformation. In particular, the intermediate reference
system defines differently:
• the X-axis of the celestial reference system; and
• the x-axis of the terrestrial reference system.
These definitions are given below, in accordance with Guinot [62], Lambert and
Bizouard [59], and McCarthy and Capitaine [60].
Resolution B1.8 adopted by the XXIV General Assembly of the International
Astronomical Union (Manchester, August 2000) recommends the use of the
“non-rotating origin” (Guinot [61]) both in the GCRS and in the ITRS. These
3.11 The Co-ordinate Transformation, Based on the Non-rotating … 387
origins are designated, respectively, as the Celestial Intermediate Origin (CIO) and
the Terrestrial Intermediate Origin (TIO). The Earth rotation angle (ERA, denoted
by h) is defined as the angle, measured along the equator, of the Celestial
Intermediate Pole (CIP) between the CIO and the TIO. Resolution B1.8 also rec-
ommends that UT1 be linearly proportional to the ERA and that the transformation
between the ITRS and GCRS be specified by the position of the CIP in the GCRS,
the position of the CIP in the ITRS, and the ERA.
With reference to the following figure, let us consider first the terrestrial refer-
ence system (ITRS). The non-rotating origin (NRO) is realised by taking into
account the motion of the CIP in the terrestrial reference system because of polar
motion. The NRO is defined by the kinematic condition of non-rotation of this point
around the rotation axis when the CIP moves in the terrestrial reference system.
Likewise, let us consider now the celestial reference system. The non-rotating
origin (NRO) is realised by taking into account the motion of the CIP in the celestial
reference system (GCRS). The NRO is defined by the kinematic condition of
non-rotation of this point around the rotation axis when the CIP moves in the
celestial reference system because of precession-nutation.
In the preceding figure, the NRO of the terrestrial reference system is the point -
on the moving equator of date, which is called the TEO. Likewise, the NRO of the
celestial reference system (which will be shown in a successive figure) is the point r
on the moving equator of date, which is called the CEO.
With reference to the celestial reference system, Kaplan [62] has proposed two
further definitions of a non-rotating origin. They are given below.
• Let a pole P(t) and a point of origin r be given on the equator. If O is the
geocentre, then the orientation of the Cartesian axes OP, Or, Oη (where Oη is
orthogonal to both OP and Or) is determined by the condition that in any
infinitesimal displacement of P there is no instantaneous rotation around
OP. Then r (also η) is a non-rotating origin on the moving equator.
• A non-rotating origin, r, is a point on the moving equator whose instantaneous
motion is always orthogonal to the equator.
These definitions are illustrated by Kaplan by means of the following figure,

which shows the motion of a non-rotating origin, r, compared to that of the true
equinox, ♈.
The green dashed line represents the trajectory described by the Celestial
Intermediate Origin, r. By definition, this trajectory is orthogonal to the moving
equator. Kaplan [31] also notes that the ecliptic, which is shown in the figure as
fixed, has actually a small motion in inertial space.
In contrast with the Celestial Intermediate Origin, r, the true equinox, ♈, has a
component of its motion along the instantaneous equator, and therefore rotates about
the Celestial Intermediate Pole (CIP). According to Capitaine and Wallace [63], the
CIO is at present very close to GCRS longitude zero and almost stationary in longitude.
By contrast, the former zero point of right ascension (i.e. the vernal equinox)
moves at over 50 arcseconds per year in GCRS longitude. Likewise, the Terrestrial
Intermediate Origin (TIO) is at present close to ITRS longitude zero and again is
almost stationary in longitude.
The Earth rotation angle (ERA), denoted by h, can be considered as the
CIO-based equivalent of the Greenwich sidereal time. The ERA is a linear function
of UT1, as follows
hðTU Þ ¼ 2pð0:7790572732640 þ 1:00273781191135448TU Þ
where TU = Julian UT1 date−2451545.0, due to the non-rotating nature of the CIO.
According to Kaplan [31], the transformation (from GCRS to ITRS) based on the
non-rotating origins can be expressed as follows
rGCRS ¼ QðtÞR3 ðhÞWðtÞrITRS
where, again, rGCRS is the position vector of the object of interest in the
Geocentric ICRS, rITRS is the corresponding vector in the International Terrestrial
Reference System, W(t) is the wobble (also called polar-motion) matrix (which
differs from the wobble matrix used in the transformation based on the equinox, as
will be shown below), R3(−h) is the elementary rotation matrix (see Sect. 3.4)
which depends on the Earth rotation angle (h) as follows
2 3
cos h sin h 0
R3 ðhÞ 4 sin h cos h 05
0 0 1
and Q(t), also denoted by [C(t)]T in Ref. [31], is the matrix which represents the
combined effects of frame bias and precession–nutation and defines the orientation
of the CIP and a longitude origin [63].
The matrix R3(−h) performs here the same function as that of the Earth rotation
matrix S(t) = R3(−hG), which has been described in Sect. 3.6 and is used in the
co-ordinate transformation based on the equinox. Instead of the Greenwich apparent
sidereal time (hG), we use here the Earth rotation angle (h).
The wobble matrix W(t) to be used here has the following expression
WðtÞ ¼ R3 ðs0 ÞR2 ðxP ÞR1 ðyP Þ
[31], where R1(a), R2(b), and R3(c) are the three elementary matrices which have
been described in Sect. 3.4, xP and yP are the angles between the Celestial
Intermediate Pole (CIP) of date and the IERS Reference Pole (IRP), and
s0 ¼ 000 :000047 T
where T = (JDE−2451545.0)/36525 is the time, expressed in Julian centuries,

elapsed between the Julian ephemeris day (JDE) of interest and J2000.0.
The shift s′ is used here because the transformation based on the non-rotating
origins must not only re-orient the pole from the ITRS z-axis to the CIP, but also
move the origin of longitude very slightly from the ITRS x-axis to the Terrestrial
Intermediate Origin, -, in accordance with Ref. [31]. The shift s′, necessary for this
purpose, is so small that its magnitude can be expressed by the approximate formula
given above (due to Lambert and Bizouard [59]), linear in time, based on the two
main circular components of polar motion [31].
By executing the matrix product in W(t) = R3(−s′) R2(xP) R1(yP), there results

cos s0 cos xP cos s0 sin xP sin yP sin s0 cos yP cos s0 sin xP cos yP sin s0 sin yP

sin s0 cos xP sin s0 sin xP sin yP þ cos s0 cos yP sin s0 sin xP cos yP þ cos s0 sin yP

sin xP cos xP sin yP cos xP cos yP
Since the three angles s′, xP, and yP are very small, the preceding matrix W(t) can
be approximated, to the first order, as follows
2 3
1 s0 xP
Wð t Þ 4 s 0 1 yP 5
xP yP 1
where xP, yP, and s′ are expressed in radians, in accordance with Kaplan [31].
By comparing the matrix W(t) given above with the corresponding matrix of
Sect. 3.7 (transformation based on the equinox), it appears that the two matrices
coincide if the value of s′ is set to zero.
It remains to show how to form the matrix Q(t), which takes account of the
motion of the Celestial Intermediate Pole in the GCRS.
Following again Kaplan [31], let us consider the following figure.
On the celestial sphere, the instantaneous (moving) equator of the Earth inter-
sects the GCRS equator at two nodes. Let N be the ascending node of the instan-
taneous equator on the GCRS equator. The scalar quantity s(t), called the CIO
locator, is the difference between the length of the arc from N westward to the CIO
(on the instantaneous equator) and the length of the arc from N westward to the
GCRS origin of right ascension (on the GCRS equator).
In the figure, r represents the CIO, and R0 represents the right ascension origin
of the GCRS (the direction of the GCRS x-axis). The definition of s implies that
s ¼ rN R0 N
where the points R0 and R′0 are equidistant from the node N. It must be stressed that
the arcs e0, s, rN, and R′0N are to be measured along the instantaneous equator, not
along a straight line. The quantity s is the extra length of the arc on the instanta-
neous equator from N to r, that is, from the ascending node to the Celestial
Intermediate Origin. The value of s depends on time and is given by
Zt
sðtÞ ¼ f½XðtÞY 0 ðtÞ YðtÞX 0 ðtÞ=½1 þ ZðtÞgdt þ s0
t0
[31], where the prime sign (′) denotes time derivatives, and X(t), Y(t), and Z(t) are
the three components of the unit vector nGCRS(t), which indicates the position of the
Celestial Intermediate Pole (CIP) in the GCRS. Then, X(t), Y(t), and Z(t) are the
co-ordinates of the CIP in the GCRS.
In the preceding figure, e0 is the arc on the instantaneous (true) equator of date
t from the CIO (denoted by r) to the equinox (denoted by ♈). The arc e0 is called
the equation of the origins. As shown in the figure, e0 is the right ascension of the
true equinox relative to the CIO (or, minus the true right ascension of the CIO). The
equation of the origins is also the difference h − hG, where h is the Earth rotation
angle and hG is the Greenwich apparent sidereal time.
The constant of integration (s0) has been set in such a way as to ensure that the
equinox-based and CIO-based computations of the Earth rotation lead to the same
results. For this purpose, the IERS Conventions [55] set s0 = 94 las (where las
stands for microarcseconds). For practical purposes, the value of the CIO locator at
any given time T is provided by a series expansion, given in Table 5.2d, p. 59, of
the IERS Conventions [56]. For convenience of the reader, this series expansion is
given below.
X
s ¼ XY=2 þ 94 þ 3808:65 T 122:68 T 2 72574:11 T 3 þ Ck sin ak
k
þ 1:73 T sin X þ 3:57 T cosð2XÞ þ 743:52 T 2 sin X
þ 56:91 T 2 sinð2F 2D þ 2XÞ þ 9:84 T 2 sinð2F þ 2XÞ 8:85 T 2 sinð2XÞ
Argument ak Amplitude Ck
X −2640.73
2X −63.53
2F − 2D + 3X −11.75
2F − 2D + X −11.21
2F − 2D + 2X +4.57
2F + 3X −2.02
2F + X −1.98
3X +1.72
l′ + X +1.41
l′ − X +1.26
l+X +0.63
l−X +0.63
where X and Y are two of the three components (X, Y, and Z) of the CIP unit vector
in the GCRS, which are to be computed as will be shown below. Software (SOFA
tools for Earth attitude) to evaluate this series is described in Ref. [48]. There is also
a table of daily values of s in Section B of The Astronomical Almanac. Let us
consider now the GCRS and the Celestial Intermediate Reference System (CIRS).
Let n, η, and f the three Cartesian axes of the CIRS. The f-axis has the direction of
nGCRS, where nGCRS is the unit vector pointing to the Celestial Intermediate Pole
(CIP). The n-axis has the direction of rGCRS, where rGCRS is the unit vector
pointing to the Celestial Intermediate Origin (CIO). The η-axis has the direction of
the unit vector y nGCRS rGCRS. Let n1, n2, and n3 (or X, Y, and Z) be the three
components of the unit vector nGCRS in the GCRS (i.e. the three non-dimensional
co-ordinates of the CIP in the GCRS). Let r1, r2, and r3 be the three components of
the unit vector rGCRS is the GCRS (i.e. the three non-dimensional co-ordinates of
the CIO in the GCRS). Let y1, y2, and y3 be the three components of the unit vector
y nGCRS rGCRS.
Then the matrix Q(t), which expresses the transformation between the
Geocentric Celestial Reference System and the Celestial Intermediate Reference
System, is
2 3 2 3
r1 y1 n1 r1 y1 X
QðtÞ ¼ 4 r2 y2 n2 5 ¼ 4 r 2 y2 Y 5 ¼ R3 ðEÞR2 ðdÞR3 ðsÞ
r3 y3 n3 r3 y3 Z
Since the co-ordinates X(t) = sin d cos E, Y(t) = sin d sin E, and Z(t) = cos d of
the CIP in the GCRS are also the components of a unit vector, then the sum of their
squares results X 2(t) + Y 2(t) + Z 2(t) = 1, and consequently
1
Z ðtÞ ¼ 1 X 2 ðtÞ Y 2 ðtÞ 2
When X(t), Y(t), and s(t) are known, the 3 3 matrix Q(t) can be formed (the
explicit time-dependence has been dropped below for convenience) as follows
2 3
1 bX 2 bXY X
4
Q ¼ bXY 1 bY 2 Y 5R3 ðsÞ
X Y 1 bð X þ Y Þ
2 2
where b = 1/(1 + Z) = 1/[1 + (1 − X 2 − Y 2)½], and R3(s) is the elementary

rotation matrix described in Sect. 3.4.
By executing the matrix product indicated above, the matrix Q results
2 3
ð1 bX 2 Þ cos s þ bXY sin s ð1 bX 2 Þ sin s bXY cos s X
4 bXY cos s ð1 bY 2 Þ sin s bXY sin s ð1 bY 2 Þ cos s Y 5
X cos s þ Y sin s X sin s Y cos s 1 bð X þ Y Þ
2 2
2 2
in accordance with Kaplan [31]. Since X +Y is much smaller than unity, then
the following function
1
b ¼ f X2 þ Y 2 ¼ 1
1 þ ½1 ðX 2 þ Y 2 Þ2
can be expanded in a Maclaurin series, around X2 + Y2 = 0, truncated after the

first-order term. Setting for convenience x X 2 + Y 2, there results
h i1
1 1
½f ð xÞx¼0 ¼ 1 þ ð 1 xÞ 2 ¼
2
h
x¼0
i
df ð xÞ 1 11 1 1
1 11 1
¼ ð1Þ 1 þ ð1 xÞ 2 ð1 xÞ ð1Þ
2 ¼ ¼
dx x¼0 2 x¼0 42 8
hence

df ð xÞ 1 1 1 1
b ¼ f ð xÞ ¼ ½f ð xÞx¼0 þ x þ O x2 þ x ¼ þ X 2 þ Y 2
dx x¼0 2 8 2 8
The complete matrix Q(t) expresses the transformation between the Geocentric
Celestial Reference System and the Celestial Intermediate Reference System. In
addition, Q(t) is consistent with the IAU 2006/2000A precession–nutation model at
the microarcsecond level. In addition to b 1=2 þ 1=8ðX 2 þ Y 2 Þ, the quantities
necessary to determine Q(t) are s, X and Y.
In order to compute the non-dimensional co-ordinates X and Y of the Celestial
Intermediate Pole as a function of time, Capitaine and Wallace [64] give the fol-
lowing expressions
X ¼ X0 þ X1 T þ X2 T 2 þ X3 T 3 þ X4 T 4 þ X5 T 5

þ Ri Rj¼0;3 as;j i t j sinðARGÞ þ ac;j i t j cosðARGÞ þ
Y ¼ Y0 þ Y1 T þ Y2 T 2 þ Y3 T 3 þ Y4 T 4 þ Y5 T 5

þ Ri Rj¼0;3 bc;j i t j cosðARGÞ þ bs;j i t j sinðARGÞ þ
where ARG stands for various combinations of the fundamentals arguments of the
nutation theory (see Sect. 3.5), including both luni-solar and planetary terms, and
T = (JDE−2451545.0)/36525 is the time (expressed in Julian centuries) elapsed
between the Julian ephemeris day (JDE) of interest and J2000.0. Capitaine and
Wallace [64] point out that, for practical reasons, the numerical expressions for
X and Y are usually multiplied by the factor 360° 60 60/(2p) = 1296000″/(2p)
in order to represent the approximate values in arcseconds of the corresponding
“angles” (strictly speaking, of their sines) with respect to the Z-axis of the GCRS. In
addition, the polynomial part of the X and Y co-ordinates originate from precession,
except for the contribution coming from the frame bias and from cross nutation
terms [64].
In practice, as has been shown above, the CIO locator (s) and the components
(X and Y) of the CIP unit vector in the GCRS can be computed by means of the
SOFA subroutines S06 and XY06, respectively [56].
The IAU 2000/2006 precession–nutation model has a high-precision level, of the
order of magnitude of microarcseconds (las). In order to reach this degree of
precision, several thousand coefficients, whose values are sometimes smaller than
1 las, are used to compute s, X, and Y.
In many cases, and in particular for satellite orbit predictions, the full model is
not necessary. For these cases, concise models have been formulated. They are fully
described in Ref. [63], from which the following table is also taken.
The table shown below compares three concise models (CPNb, CPNc, and
CPNd) with the full model (Reference) and also with the IAU 2000B model.
The coefficients used in the three concise models are also given in three
respective tables contained in Ref. [63].
Model Coefficients Frequencies (mas) RMS (mas) Worst Speed

Reference 4006 1309 – – 1
IAU 2000B 354 77 0.28 0.99 7.6
CPNb 227 90 0.28 0.99 15.3
CPNc 45 18 5.4 16.2 138
CPNd 6 2 160 380 890
A short account of the three models is given below. The first of them, known as
the CPNb model, has an accuracy of 1 mas in the time span 1995–2050. Its per-
formance is about that of the IAU 2000B model. As shown in the table, the number
of coefficients used in CPNb is 227 (220 for X and Y plus 7 for s + XY/2) plus those
required for the fundamental arguments.
In the CPNb model, the approximation in the X and Y series is of the order of:
• the 4th order in the precession quantity w;
• the 2nd order in the nutation quantities Dw and De;
• the 1st order in the precession quantities (x–e0J2000), (e0D–e0J2000) and v;
• the 1st order in the frame-bias quantities n0, η0 and da0; and
• the 3rd order in the cross terms between the precession quantity w and nutation
(in other words, there are terms of the form w2 nutation).
The X and Y series are expressed as follows

X ¼ n0 þ w sin e0J2000 w3 =6 sin e0J2000 þ wðx e0J2000 Þ cos e0J2000 þ Dw sin e0J2000
þ DwDe cos e0J2000 þ ðw cos e0J2000 vÞDe þ ðe0D e0J2000 ÞDw cos e0J2000

w2 =2 Dw sin e0J2000

Y ¼ g0 þ ðx e0J2000 Þ w2 =2 sin e0J2000 cos e0J2000 þ w4 =24 sin e0J2000 cos e0J2000

þ da0 w sin e0J2000 þ De Dw2 =2 sin e0J2000 cos e0J2000

ðw cos e0J2000 vÞDw sin e0J2000 w2 =2 cos e20J2000 De
Likewise, the approximation of CPNb in the s + XY/2 series is of the order of:
• the 1st order in the largest polynomial development coefficients (Xi, Yi)i=0,2 of
X and Y; and
• the 1st order in the coefficients ai and bi of the sine and cosine Fourier terms of
arguments (xiT − ui) in X and Y, respectively.
The s + XY/2 series is expressed as follows
" #
X
s þ XY=2 ¼ X1 Y0 þ 1=2 ðai bi xi Þ T þ 1=3X1 Y2 T 3 þ ð1=4ai bi Þ sin½2ðxi T ui Þ
i
þ ðbi =xi ÞX1 sinðxi T ui Þ þ Y2 T 2 ai sinðxi T ui Þ
The coefficients for the CPNb model are given in Appendix E to Ref. [63].
The second model, known as the CPNc model, has an accuracy of 16 mas in the
same time span as that indicated above. The number of coefficients used in CPNc is
45 (42 for X and Y plus 3 for s + XY/2) plus 3 for the fundamental arguments.
In the CPNc model, all coefficients of T 2 and above are removed from the X and
Y series. A three-coefficient series is given for s + XY/2. The three expressions
given above (for, respectively, X, Y, and s + XY/2) reduce to those shown below.

X ¼ n0 þ w sin e0J2000 w3 =6 sin e0J2000 þ Dw sin e0J2000
þ wDe cos e0J2000 þ ðe0D e0J2000 ÞDw cos e0J2000

Y ¼ g0 þ ðx e0J2000 Þ w2 =2 sin e0J2000 cos e0J2000
þ De wDw sin e0J2000 cos e0J2000
s þ XY=2 ¼ 1=2Ri ðai bi xi ÞT þ 1=3X1 Y2 T 3 þ ðbi =xi ÞX1 sinðxi T ui Þ
The coefficients for the CPNc model (see Appendix F to Ref. [63]) are given below.
Term Amplitude (mas) l l′ F D X

X −17251
X T 2004191898
X T2 −429783
X T3 −198618
Y −5530
Y T −25896
Y T2 −22407275
s + XY/2 T 3809
s + XY/2 T3 −72574
X sin −6844318 0 0 0 0 1
X T sin −3310 ′′ ′′ ′′ ′′ ′′
X T cos +205833 ′′ ′′ ′′ ′′ ′′
Y cos +9205236 ′′ ′′ ′′ ′′ ′′
Y T sin +153042 ′′ ′′ ′′ ′′ ′′
s + XY/2 sin −2641 ′′ ′′ ′′ ′′ ′′
(continued)
(continued)
X sin +82169 0 0 0 0 2
Y cos −89618 ′′ ′′ ′′ ′′ ′′
X sin +2521 0 0 0 2 0
X sin +5096 0 0 2 −2 1
Y cos −6918 ′′ ′′ ′′ ′′ ′′
X sin −523908 0 0 2 −2 2
X T cos +12814 ′′ ′′ ′′ ′′ ′′
Y cos +573033 ′′ ′′ ′′ ′′ ′′
Y T sin +11714 ′′ ′′ ′′ ′′ ′′
X sin −15407 0 0 2 0 1
Y cos +20070 ′′ ′′ ′′ ′′ ′′
X sin −90552 0 0 2 0 2
Y cos +97847 ′′ ′′ ′′ ′′ ′′
X sin −8585 0 1 −2 2 −2
Y cos −9593 ′′ ′′ ′′ ′′ ′′
X sin +58707 0 1 0 0 0
Y cos +7387 ′′ ′′ ′′ ′′ ′′
X sin −20558 0 1 2 −2 2
Y cos +22438 ′′ ′′ ′′ ′′ ′′
Y cos +2555 1 0 −2 −2 −2
X sin −4911 1 0 −2 0 −2
Y cos −5331 ′′ ′′ ′′ ′′ ′′
X sin −6245 1 0 0 −2 0
Y cos +3144 1 0 0 0 −1
X sin +28288 1 0 0 0 0
X sin +2512 1 0 0 0 1
Y cos −3324 ′′ ′′ ′′ ′′ ′′
Y cos +2636 1 0 2 0 1
X sin −11992 1 0 2 0 2
Y cos +12903 ′′ ′′ ′′ ′′ ′′
The third model, known as the CPNd model, has an accuracy of 390 mas in the
same time span. The number of coefficients used in CPNd is 6 (4 for X and Y plus 1
for a T-term in X plus 1 for a T2-term in Y). There is no term for s + XY/2, because
s is set to zero.
In the CPNd model, the X and Y series include only four coefficients, which are a
T-term in X, a T2-term in Y, and the main 18.6-year nutation term. The three
expressions given above (for, respectively, X, Y, and s + XY/2) reduce to those
shown below.
X ¼ w sin e0J2000 þ Dw sin e0J2000

Y ¼ w2 =2 sin e0J2000 cos e0J2000 þ De
s¼0
The coefficients for the CPNd model (see Appendix G to Ref. [63]) are given below.

X T 2004191898
Y T2 −22407275
X sin −6844318 0 0 0 0 1
Y cos +9205236 ′′ ′′ ′′ ′′ ′′
X sin −523908 0 0 2 −2 2
Y cos +573033 ′′ ′′ ′′ ′′ ′′
By using the three concise models indicated above, the computation costs are
reduced by 1, 2, and 3 orders of magnitude, respectively, in comparison with the
full model.
3.12 The Co-ordinate Transformation, According

to the GOCE Standards, Between the Celestial
and Terrestrial Reference Systems
The Gravity field and steady-state Ocean Circulation Explorer (GOCE) is a satellite
launched on the 17th of March 2009 by the European Space Agency (ESA) and
meant to measure the Earth gravity field by means of accelerometers and provide a
model of the geoid with high accuracy (of the order of magnitude of 1 10−5 m/s2)
and spatial resolution (better than 100 km).
The co-ordinate transformation described by Gruber et al. [65], according to the
GOCE standards, is substantially the classical transformation based on the equinox
and the Greenwich apparent sidereal time. Using the same notation as that of the
previous sections, this transformation can be expressed in the form
rICRS ¼ BðtÞPðtÞNðtÞSðtÞWðtÞrITRS
where B(t) is the frame-bias matrix, P(t) is the precession matrix, N(t) is the
nutation matrix, S(t) is the Earth rotation matrix, and W(t) is the polar-motion
matrix.
Since these matrices differ from the homonymous matrices which have been
shown in the preceding sections, they are briefly described below.
The frame-bias matrix B(t) is defined in the GOCE standards as follows
BðtÞ ¼ R3 ðda0 ÞR2 ðdwB sin e0J2000 dX ÞR1 ðdeB þ dY Þ
where R1(a), R2(b), and R3(c) are the elementary rotation matrices defined in
Sect. 3.4, e0J2000 = 84381″.448 is the obliquity of the ecliptic at J2000.0 with
respect to the mean equator at J2000.0, dwB = −0″.041775 and deB = −0″.0068192
are, respectively, the longitude and obliquity of the celestial pole, da0 = −0″.0146
is the offset in the ICRS right ascension origin with respect to the dynamical
equinox of J2000.0, as measured in an inertial (non-rotating) reference system, and
dX and dY are the celestial pole offsets determined from VLBI observations.
Daily values of these offsets are published by the IERS within the EOP (IERS)
08 C04 series [66]. Gruber et al. [65] take e0J2000 = 84381″.448, which is the old
value given by Lieske et al. [67] instead of the new value 84381″.406 used by
Kaplan [31]. Gruber et al. [65, p. 17, first line] point out that the values of the
celestial pole offsets must be interpolated to the epoch of computation. This makes
the frame-bias matrix B(t) of [65] a function of time, because of the presence of dX
and dY in the definition of B(t).
The arguments b and a of the elementary rotation matrices R2(b) and R1(a)
contained in the formula written above are expressed (in milliarcseconds, mas) in
the IERS Conventions–2003 [55] as follows
b ¼ 16:617 0:0016 T 2 þ 0:0007 cos X

a ¼ 6:8192 0:1419 T þ 0:0005 sin X
where X is the longitude of the mean ascending node of the Moon orbit on the
ecliptic, measured from the mean equinox of date, and T is the time, expressed in
Julian centuries, elapsed between the Julian ephemeris day of interest and J2000.0,
that is, T = (JDE−2451545.0)/36525.
These expressions show that the values of the celestial pole offsets are of the
order of magnitude of milliarcseconds, whereas the corrective terms to be added to
dwB sin e0J2000 = −16.617 mas and to deB = −6.8192 mas are smaller by one or
two orders of magnitude. By executing the matrix product indicated above, the
resulting matrix B(t) is equal to the transpose of the constant matrix B of Sect. 3.9,
if dX and dY are set to zero.
The precession matrix P(t) is computed by Gruber et al. [65] by means of the
canonical 4-rotation method, as follows
PðtÞ ¼ R1 ðe0J2000 ÞR3 ðwÞR1 ðxÞR3 ðvÞ
where, again, e0J2000 = 84381″.448, and the angles w, x, and v are
w ¼ 503800 :47875 T 100 :07259 T 2 000 :001147 T 3

x ¼ e0J2000 000 :02524 T þ 000 :05127 T 2 000 :007726 T 3
v ¼ 1000 :55260 T 200 :38064 T 2 000 :001125 T 3
3.12 The Co-ordinate Transformation, According to the GOCE … 399
which are the same expressions as those given by Lieske et al. [67]. By contrast,
Kaplan [31] gives the new expressions of the angles w, x, and v, which have been
shown in Sect. 3.4. Again, T = (JDE−2451545.0)/36525 is the time (expressed in
Julian centuries) elapsed between the Julian ephemeris day of interest and J2000.0.
JDE is the Julian day (JD) corresponding to an instant of time measured in the
Terrestrial Timescale.
The nutation matrix N(t) is defined by Gruber et al. as follows
NðtÞ ¼ R1 ðe0D ÞR3 ðDwÞR1 ðe0D þ DeÞ
where e0D (the mean obliquity of the ecliptic of date with respect to the mean
equator of date) depends on time as follows [65]:
e0D ¼ 8438100 :448 4600 :84024 T 000 :00059 T 2 þ 000 :001813 T 3
Dw and De are the nutations in, respectively, longitude and obliquity, which are
due to the sums of luni-solar (LS) and planetary (PL) terms, as will be shown
below. The luni-solar terms (DwLS and DeLS) of the nutation angles are given by the
following trigonometric series
X
662
DwLS ðTÞ ¼ ½ASk sin ak ðTÞ þ ACk cos ak ðTÞ
k¼1
X
38
þT A0Sk sin ak ðTÞ þ A0Ck cos ak ðTÞ
k¼1
X
662
DLS ðTÞ ¼ ½BSk sin ak ðTÞ þ BCk cos ak ðTÞ
k¼1
X
38
þT B0Sk sin ak ðTÞ þ B0Ck cos ak ðTÞ
k¼1
where the arguments ak(T) of the trigonometric functions are linear combinations
X
5
ak ðT Þ ¼ nkj F j ðT Þ
j¼1
of the fundamental arguments F1(T), F2(T), …, F5(T) of the luni-solar nutation. The
argument multipliers nkj have integral values, which are given in Attachment 11.4
to the GOCE Standards [65]. The values of the coefficients ASk, ACk, A′Sk, A′Ck, BSk,
BCk, B′Sk, and B′Ck are also given in Attachment 11.4 cited above.
The fundamental arguments F1(T), F2(T), …, F5(T) are the Delaunay variables,
which are expressed by the following polynomials.
Mean anomaly of the Moon:
F 1 ðTÞ l ¼ 48586800 :249036 þ 171791592300 :2178 T þ 3100 :8792 T 2

þ 000 :051635 T 3 000 :00024470 T 4
Mean anomaly of the Sun:
F 2 ðTÞ l0 ¼ 128710400 :79305 þ 12959658100 :0481 T 000 :5532 T 2

þ 000 :000136 T 3 000 :00001149 T 4
Mean longitude of the Moon (L) minus mean longitude of the Moon ascending
node (X):
F 3 ðTÞ F ¼ 33577900 :526232 þ 173952726200 :8478 T 1200 :7512 T 2

000 :001037 T 3 þ 000 :00000417 T 4
Mean elongation of the Moon from the Sun:
F 4 ðTÞ D ¼ 107226000 :70369 þ 160296160100 :2090 T 600 :3706 T 2

þ 000 :006593 T 3 000 :00003169 T 4
Mean longitude of the ascending node of the Moon:
F 5 ðTÞ X ¼ 45016000 :398036 696289000 :5431 T þ 700 :4722 T 2 þ 000 :007702 T 3

000 :00005939 T 4
The polynomials given above differ from those of Sect. 3.5 (which, in turn, are
those of the 1980 theory of nutation, taken from Seidelmann [42]). In addition, as
has been shown above, the trigonometric series expressing the nutation angles Dw
and De are also different. The planetary terms (DwPL and DePL) of the nutation
angles are given by the following trigonometric series
X
658
DwPL ðT Þ ¼ ½CSk sin bk ðT Þ þ CCk cos bk ðT Þ
k¼1
X
658
DPL ðT Þ ¼ ½DSk sin bk ðT Þ þ DCk cos bk ðT Þ
k¼1
where the arguments bk(T) of the trigonometric functions are linear combinations
X
14
bk ðT Þ ¼ nkj f j ðT Þ
j¼1
of the fundamental arguments f1(T), f2(T), …, f14(T) of the planetary nutation.

The argument multipliers nkj are integers, whose values are given in Attachment
11.5 to the GOCE Standards [65]. The values of coefficients CSk, CCk, DSk, and DCk
are also given in Attachment 11.5 cited above.
The first five of the fourteen fundamental arguments f1(T), f2(T), …, f14(T) are
the same as the five fundamental arguments F1(T), F2(T), …, F5(T) of the luni-solar
terms. However, simplified expressions of F1(T), F2(T), …, F5(T) are used by
Gruber et al. [65] to compute the planetary terms of nutation. Therefore, the
complete set of f1(T), f2(T), …, f14(T) is given below.
f 1 ðTÞ l ¼ 2:355555980 þ 8328:6914269554 T

f 2 ðTÞ l0 ¼ 6:240060130 þ 628:3019550000 T
f 3 ðTÞ F ¼ 1:627905234 þ 8433:4661581310 T
f 4 ðTÞ D ¼ 5:198466741 þ 7771:3771468121 T
f 5 ðTÞ X ¼ 2:182439200 33:7570450000 T
f 6 ðTÞ lMe ¼ 4:402608842 þ 2608:7903141574 T
f 7 ðTÞ lVe ¼ 3:176146697 þ 1021:3285546211 T
f 8 ðTÞ lE ¼ 1:753470314: þ 628:3075849991 T
f 9 ðTÞ lMa ¼ 6:203480913 þ 334:0612426700 T
f 10 ðTÞ lJu ¼ 0:599546497 þ 52:9690962641 T
f 11 ðTÞ lSa ¼ 0:874016757 þ 21:3299104960 T
f 12 ðTÞ lUr ¼ 5:481293871 þ 7:4781598567 T
f 13 ðTÞ lNe ¼ 5:321159000 þ 3:8127774000 T
f 14 ðTÞ pa ¼ 0:0243817500 T þ 0:00000538691 T 2
The coefficients of these polynomials are expressed in radians, and T is the time,
expressed in Julian centuries, elapsed between the Julian ephemeris day of interest
and J2000.0, that is, T = (JDE−2451545.0)/36525.
The Earth rotation matrix S(t) is
SðtÞ ¼ R3 ðhG Þ
where hG is the Greenwich apparent sidereal time. The latter, in turn, is computed
by Gruber et al. [65] as follows:
hG ¼ hðTU Þ þ pðTÞ þ Dw cos e0D þ qðTÞ
where h(TU) is the Earth rotation angle, Dw is the nutation in longitude (see
Sect. 3.5), TU = Julian UT1 date −2451545.0 is the number of days elapsed since
J2000.0, e0D = 84381″.448 − 46″.84024 T − 0″.00059 T2 + 0″.001813 T3 is the
mean obliquity of the ecliptic of date with respect to the mean equator of date, T is
the time, expressed in Julian centuries, elapsed between the Julian ephemeris day of
interest and J2000.0, and p(T) and q(T) are two functions of time T, as will be
shown below.
First of all, the sum h(TU) + p(T) is the Greenwich mean sidereal time hG0 (see
Sect. 3.6), and consequently the preceding expression of the Earth rotation angle
hG ¼ hG0 þ Dw cos e0D þ qðTÞ
Then, as has been shown in Sects. 3.6 and 3.11, the Earth rotation angle, h(TU),
is a first-degree polynomial in TU = Julian UT1 date−2451545.0, as follows
hðTU Þ ¼ 2pð0:7790572732640 þ 1:00273781191135448 TU Þ
Remembering the definitions given in Chap. 2, Sect. 2.4, there results
UT1 ¼ UTC þ ðUT1UTCÞ ¼ TAI þ ðUTCTAIÞ + (UT1 UTC)

¼ TT þ ðTAITTÞ þ ðUTCTAIÞ þ ðUT1UTCÞ
TAI TT ¼ 32:184 s
Information on the difference UTC–TAI is published by the IERS in the

Bulletin C [68]. The instantaneous value of the difference UT1–UTC can be
obtained, after interpolating the daily values published by the IERS within the EOP
(IERS) 08 C04 series [66], by applying an additional component, denoted
DUT1tidal(T), which takes account of the effect of the ocean tides with periods less
than two days, as follows
ðUT1 UTCÞT ¼ ðUT1 UTCÞIERS

T þ DUT1tidal ðTÞ
Taking account of this additional component and defining
TU0 ¼ JDðTTÞ 2451545:0

h i
þ ðUT1 UTCÞIERS
T þ ðTAI TTÞ þ ðUTC TAIÞ =86400
the argument TU of the Earth rotation angle h(TU) can be written as follows
TU ¼ TU0 þ DUT1tidal ðTÞ=86400
The additional component DUT1tidal(T), due to the ocean tides, is made up by 41

diurnal and 30 semi-diurnal constituents, as follows
X
71
DUT1tidal ðT Þ ¼ tidal
USk sin nk ðT Þ þ UCk
tidal
cos nk ðT Þ
k¼1
The arguments nk(T) of the sine and cosine functions, in turn, result from linear
combinations
X
5
nk ð T Þ ¼ nkj F j ðT Þ þ nk6 vðtu Þ
j¼1
The argument multipliers nkj are integers, and the coefficients Utidal tidal
Sk and UCk are
real numbers, whose values are given in Attachment 11.1 of [65]. They are also
given in Table 8.3 of the IERS Technical Note No. 32 [55].
The fundamental arguments F1(T) l, F2(T) l′, …, F5(T) X are the
fourth-degree polynomials given above for the nutation matrix.
According to Gruber et al. [65], the quantity v(tu) in the last term of the
expression written above is GMST + p, where GMST stands for Greenwich mean
sidereal time and is computed [49, p. 360, Eq. 14] as follows

SðT 0 Þ ¼ 67310s :54841 þ 876600h þ 8640184s :812866 T 0 þ 0s :093104 T 02
0s :0000062 T 03
where T′ is the number of Julian centuries (consisting of 36525 days of 86400 s of

dynamical time each) elapsed since JD 2451545.0 TDB, and S is dynamical sidereal
time. The quantity v(tu) results from
vðtu Þ ¼ p þ GMST ¼ 12h þ SðT 0 Þ ¼ ð12 3600Þs þ 67310s :54841

þ ½ð876600 3600Þs þ 8640184s :812866tu þ 0s :093104 tu2 0s :0000062 tu3
where
tu ¼ TU0 =36525 ¼ fJDðTTÞ 2451545:0 þ ½ðUT1 UTCÞIERS

T þ ðTAI TTÞ
þ ðUTC TAIÞ=86400g=36525
The polynomial p(T) appearing in hG = h(TU) + p(T) + Dw cos e0D + q(T) is the
polynomial of Terrestrial Time, that is,
pðTÞ ¼ 000 :014506 þ 461200 :15739966 T þ 100 :39667721 T 2 000 :00009344 T 3

þ 000 :00001882 T 4
where T = (JDE−2451545.0)/36525 is the time (expressed in Julian centuries)

elapsed between the Julian ephemeris day of interest and J2000.0.
The term q(T) appearing in hG = h(TU) + p(T) + Dw cos e0D + q(T) is the
non-polynomial part of the Greenwich apparent sidereal time. This term is
expressed by the following trigonometric series
X
33 0
0
qð T Þ ¼ ½ESk sin ck ðT Þ þ ECk cos ck ðT Þ þ T ES1 sin c1 ðT Þ þ EC1 cos c1 ðT Þ
k¼1
where, again, T = (JDE−2451545.0)/36525 is the time (expressed in Julian cen-

turies) elapsed between the Julian ephemeris day of interest and J2000.0.
The arguments ck(T) of the sine and cosine functions are linear combinations
X
5 X
14
ck ð T Þ ¼ nkj F j ðT Þ þ nkj f j ðT Þ
j¼1 j¼6
of the fundamental arguments of the luni-solar and planetary nutation, the argument
multipliers nkj have integral values, and the coefficients ECk, ESk, E′S1, and E′C1 have
real values, all of which are given in Attachment 11.8 to Ref. [65].
As has been shown above, the computation of instantaneous values of the dif-
ference UT1–UTC in the expression
ðUT1 UTC) ¼ ðUT1 UTCÞIERS þ DUT1tidal ðT Þ
requires an interpolation of the daily values published by the IERS within the EOP
(IERS) 08 C04 series [66]. For consistency with the techniques applied to derive
these combined series, Gruber et al. [65] first subtract the long-periodic zonal tides,
dUT1zonal(Tnode), from the values (UT1—UTC)IERS node at the interpolation nodes, as
follows
ðUT1Red UTCÞIERS IERS

node ¼ ðUT1 UTCÞnode dUT1zonal ðT node Þ
and then add the zonal tides to the interpolated values, (UT1Red – UTC)IERS
T , as
follows
ðUT1 UTCÞIERS
T ¼ ðUT1Red UTCÞIERS
T þ dUT1zonal ðTÞ
By so doing, the preceding expression
ðUT1 UTCÞT ¼ ðUT1 UTCÞIERS

T þ DUT1tidal ðTÞ
becomes
ðUT1 UTCÞT ¼ ðUT1Red UTCÞIERS

T þ dUT1zonal ðTÞ þ DUT1tidal ðTÞ
The zonal variation, dUT1zonal(T), includes 62 components of periods longer

than 5 days, in the following form
X
62
dUT1zonal ðT Þ ¼ zonal
USk sin gk ðT Þ þ UCk
zonal
cos gk ðT Þ
k¼1
where the arguments ηk(T) of the trigonometric functions are linear combinations
X
5
gk ðT Þ ¼ nkj F j ðT Þ
j¼1
of the fundamental arguments F1(T), F2(T), …, F5(T) of the luni-solar nutation. The
argument multipliers nkj have integral values, and the coefficients Uzonal
Sk and Uzonal
Ck
have real values, all of which are given in Attachment 11.9 to Ref. [65] or in
Table 8.1 of Ref. [55].
The polar-motion matrix W(t) is defined by Gruber et al. [65] as follows
WðtÞ ¼ R3 ðs0 ÞR2 ðxP ÞR1 ðyP Þ
where s′ is the position of the Terrestrial Intermediate Origin (TIO) on the equator
of the Celestial Intermediate Pole (CIP), both of which are described in Sect. 3.11.
Gruber et al. compute s′ as indicated by Lambert and Bizouard [59]:
s0 ¼ 000 :000047 T
where T = (JDE−2451545.0)/36525 is the time, expressed in Julian centuries,

elapsed between the Julian ephemeris day (JDE) of interest and J2000.0. The angles
xP and yP are the co-ordinates of the Celestial Intermediate Pole (CIP) of date with
respect to the IERS Reference Pole (IRP).
Gruber et al. [65] compute the instantaneous values of xP and yP as follows.
First, the EOP (IERS) 08 C04 daily series [66] are interpolated for the time (T) of
interest. Then, two terms are added to the interpolated values to take account of the
ocean tides and the nutation terms with periods less than two days, as shown below:

xP x Dxtidal ðT Þ Dxnutation ðT Þ
¼ IERS þ þ
yP T
yIERS T
Dytidal ðT Þ Dynutation ðT Þ
The tidal corrections, Dxtidal(T) and Dytidal(T), include 41 diurnal and 30

semi-diurnal terms, as follows
X
71
Dxtidal ðT Þ ¼ tidal
XSk sin nk ðT Þ þ XCk
tidal
cos nk ðT Þ
k¼1
X
71
Dytidal ðT Þ ¼ tidal
YSk sin nk ðT Þ þ YCk
tidal
cos nk ðT Þ
k¼1
The tidal terms indicated above are exactly the same as those of the preceding
expression
X
71
DUT1tidal ðT Þ ¼ tidal
USk sin nk ðT Þ þ UCk
tidal
cos nk ðT Þ
k¼1
Consequently, the argument nk(T) of the sine and cosine functions results from
the preceding expression
X
5
nk ð T Þ ¼ nkj F j ðT Þ þ nk6 vðtu Þ
j¼1
The numerical values of the coefficients Xtidal tidal tidal,

Sk , XCk , YSk and Ytidal
Ck are given in
Attachment 11.1 to Ref. [65] or in Table 8.2 of Ref. [55].
The nutational corrections, Dxnutational(T) and Dynutational(T), to the CIP
co-ordinates include 10 diurnal and sub-diurnal terms, as follows
X
10
Dxnutational ðT Þ ¼ nutational
XSk sin fk ðT Þ þ XCk
nutational
cos fk ðT Þ
k¼1
X
10
Dynutational ðT Þ ¼ nutational
YSk sin fk ðT Þ þ YCk
nutational
cos fk ðT Þ
k¼1
where the arguments fk(T) of the trigonometric functions are linear combinations
X
5
fk ð T Þ ¼ nkj F j ðT Þ þ nk6 vðtu Þ
j¼1
of the fundamental arguments of the luni-solar nutation. The argument multipliers

nkj are integers, and the coefficients Xnutational
Sk , Xnutational
Ck , Ynutational,
Sk and Ynutational
Ck are
real numbers, whose values are given in Attachment 11.6 of Ref. [65]. They are
also given in Table 5.1, second part, of Ref. [55].
Gruber et al. [65] have compared the results obtained by them by using the
GOCE standards in some test cases with the results coming from the application of
the IERS transformation matrices, which have been described in the preceding
sections. According to Gruber et al., the deviations observed do not exceed
4.5 10−12 [65].
3.13 The Luni-Solar Perturbation 407
3.13 The Luni-Solar Perturbation
The Earth has a natural satellite, the Moon, and revolves about the Sun. The
gravitational forces exerted by these bodies on Earth artificial satellites are not
negligible, especially for those of them which orbit the Earth at high altitudes.
The presence of other bodies in the gravitational field exerted by a main central
body makes the problem of orbit determination a many-body problem. The general
n-body problem consists in determining the motion of an isolated set of n bodies,
having masses m1, m2, …, mn, attracting one another with Newtonian gravitational
forces. In other words, given the positions and velocities of the n bodies, considered
as point-masses, at an initial time t0, we want to determine the positions and
velocities of such bodies at any subsequent time t. As is well known (see, i.e. [13]),
the general n-body problem cannot be solved analytically.
In case of artificial satellites orbiting the Earth at low altitudes, the gravitational
force due to the Earth is by far stronger than those exerted by the Moon and the
Sun. Therefore, the problem can be solved by using the methods of the perturbation
theory, as will be shown below. In particular, since the masses of artificial satellites
are negligible in comparison with those of natural celestial bodies, then the masses
of such satellites do not perturb appreciably the gravitational field which acts upon
them. By inspection of the table (due to Blitzer [1]) shown in Sect. 3.1, the
luni-solar acceleration appears to have the constant value of 10−6 m/s2 at orbital
altitudes below 1500 km, whereas the main term of the Earth gravitational accel-
eration decreases with altitude from 9.35 to 6.42 m/s2. Things change in case of
geostationary satellites (h = 35786 km), which are subject to a luni-solar acceler-
ation of 7 10−6 m/s2 and to a central gravitation acceleration of 0.22 m/s2. As a
result of the luni-solar acceleration, together with those coming from the
non-spherical Earth and the solar radiation pressure, the orbital plane of a geosta-
tionary satellite is subject to change. Since the mean obliquity at J2000.0 of the
ecliptic plane is about 23°.44 and the orbital plane of the Moon is inclined of 5°.14
with respect to the ecliptic, then the secular terms of the rate of change of incli-
nation i′sec of the orbit of a geostationary satellite with respect to the equator are
respectively
0
isec ¼ 0 :478 per annum ðfor the MoonÞ
0 M
isec S ¼ 0 :319 per annum ðfor the SunÞ
Therefore, in accordance with Walter [3], the total luni-solar rate of change of
orbital inclination with time amounts to
0
isec M þ S ¼ 0 :797 per annum
In the absence of corrective actions, called north–south station-keeping, the

plane of a satellite orbit precesses, causing an increase or decrease of orbit incli-
nation with respect to the equator. The beam of the satellite antenna draws a “figure
of eight” on the surface of the Earth [69].
Walter [3] has computed the velocity correction (Dv) necessary to a geosta-
tionary satellite in a year to correct the perturbation induced by the Moon and the
Sun, as follows
" #

i0sec MþS 0 :797
Dv ¼ 2v sin ¼ 2 3:0747 sin
2 2
¼ 0:042769 km/s per annum
or 42.77 m/s per annum. The value 3.0747 is the velocity, expressed in km/s, of a
geostationary satellite, which is a satellite revolving around the Earth in a circular
equatorial orbit whose radius, rG, is such that the period of revolution, TG, of the
satellite is equal to the period (one sidereal day, that is, 23h 56m 4s.0916, that is,
86164.0916 s) of rotation of the Earth around its axis, as seen from the fixed stars.
From the third law of Kepler (Chap. 1, Sect. 1.3), there results
4p2 3
TG2 ¼ r
lE G
where lE = 398600.4418 km3/s2 is the gravitational parameter of the Earth sub-

stituting the values of TG and lE in the preceding expression and solving for rG,
there results rG = 42164 km with five significant figures. This value, substituted in
the energy integral (v2G = lE/rG), leads to vG = 3.0747 km/s. According to
Pustynski [70], at an altitude h of about 50000 km, the luni-solar acceleration
becomes greater than the acceleration caused by the Earth oblateness.
Before showing how to compute the luni-solar perturbation, it is necessary to
study the general problem of the motion of an object revolving about a central mass
(e.g. an artificial satellite in orbit about the Earth) and also subject to the perturbing
gravitational force exerted by a third body (e.g. the Moon or the Sun). Following
Murray [71], let us consider the motion of two masses (mi and mj) orbiting a central
mass (mc), as shown in the following figure.
If we use the following notation

1
jri j ri ¼ x2i þ y2i þ z2i 2
1
r j r j ¼ x2 þ y2 þ z 2 2
j j j
h 2 i1
rj ri ¼ xj xi þ yj yi 2 þ zj zi 2 2
then the equations of motion of the three masses mc, mi, and mj, with respect to the
origin O, are respectively

!
ri rj
mc R00c ¼ Gmc mi 3 þ Gmc mj 3
ri rj
!

rj ri ri
mi R00i ¼ Gmi mj Gm m
i c
rj ri 3 ri3
! !
00 ri rj rj
mj Rj ¼ Gmj mi Gmj mc 3
ri rj 3 rj
The equations of motion of the masses mi and mj, with respect to the central mass
mc, are
r00i ¼ R00i R00c

r00i ¼ R00j R00c
Substituting

!
ri rj
R00c ¼ Gmi 3 þ Gmj 3
ri rj
!

rj ri ri
R00i ¼ Gmj Gmc
rj ri 3 ri3
! !
ri rj rj
R00j ¼ Gmi Gmc
ri rj 3 rj3
into
r00j ¼ R00i R00c

r00j ¼ R00j R00c
leads to

!
ri rj ri rj
r00i þ Gðmc þ mi Þ 3 ¼ Gmj
ri rj ri 3 r 3j
! !
rj ri rj ri
r00j þ G mc þ mj ¼ Gmi
rj3 ri rj 3 r 3i
If we set

Gðmc þ mi Þ G mc þ mj
Ui ¼ Uj ¼
ri rj
then the preceding equations can be written as follows
r00i ¼ $i ðU i þ F i Þ gradi ðU i þ F i Þ

r00j ¼ $j U j þ F j gradj U j þ F j
where gradi (or gradj) designates the gradient with respect to the co-ordinates xi, yi,
and zi (or xj, yj, and zj), and
Gmj r r
F i ¼ Gmj i 3 j
rj ri rj
Gmi r r
F j ¼ Gmi i 3 j
ri rj ri
The function Fi is called the disturbing function for the mass mi due to the
presence of the mass mj. The disturbing function is substantially the perturbing
potential of gravitation to which a body, revolving around a central body, is subject
because of the presence of a third body. Now we change the notation, so that the
orbiting masses mi and mj become, respectively, m and m*, and their position
vectors ri and rj become, respectively, r and r*. Without loss of generality, we can
assume r < r*, so that m is the inner mass and m* is the outer mass. Consequently,
the equation of motion written above for the inner mass m

!
ri rj ri rj
r00i þ Gðmc þ mi Þ 3 ¼ Gmj
ri rj ri 3 rj3
is rewritten in the new notation as follows

!
r r r r
r00 þ Gðmc þ mÞ ¼ Gm 3
r3 j r rj3 r
and the disturbing function for the inner mass is rewritten as follows
Gm rr

l rr

F¼ Gm ¼ l
jr rj r 3 jr rj r 3
where l* = Gm* is the gravitational parameter relating to the outer mass.

Likewise, the equation of motion and the disturbing function for the outer mass
m* become respectively

!
00 r r r r
r þ Gðmc þ m Þ 3 ¼ Gm
r j r r j 3 r 3
and
Gm r r l r r
F ¼ Gm ¼ l
j r r j r3 jr r j r3
where l = Gm is the gravitational parameter relating to the inner mass.

Now, operating as in Sect. 3.2, let c be the angle between the vectors r and r*, as
The following geometrical relation (cosine rule) holds

r 2 r 12
1
jr rj ¼ r 2 þ r 2 2rr cos c 2 ¼ r 1 þ 2 cos c
r r
As has been shown in Sect. 3.2, by taking the reciprocal of |r* − r| and
expanding 1/|r* − r| in a series of Legendre polynomials, there results
r 2 r 12 1 k
1 1 1X r
¼ 1 þ 2 cos c ¼ Pk ðcos cÞ
jr rj r r r r k¼0 r
where Pk(cos c) is the Legendre polynomial of kth degree. This series converges
when r < r*, according to the assumption made above.
Since P0(cos c) = 1 and P1(cos c) = cos c = (r r*)/(rr*), then the first term
(corresponding to k = 0) of the sum is 1/r*, and the second term (k = 1) of the sum
is (1/r*)(r/r*) cos c = (1/r*)(r/r*)(r r*)/(rr*) = (r r*)/r*3. According to Murray
[71], the first term, which does not depend on r, can be neglected, and the second
term, multiplied by l*, vanishes with −l*(r r*)/r*3 in
l l ðr r Þ
F¼
jr rj r 3
which is the expression of the disturbing function, F, for the inner mass. Therefore,
the expansion of F in a series of Legendre polynomials is
1 k
l X r
F¼ Pk ðcos cÞ
r k¼2 r

Murray also expresses the expansion of the disturbing function as a function of

the orbital elements of the outer, or perturbing, body (denoted by an asterisk) and
the orbital elements of the inner body (without asterisk), as follows:
X
F ¼ l Sða; a ; e; e ; i; i Þ cos u
where S is a function whose form is to be determined, and u is an angle such that
u ¼ j1 k þ j2 k þ j3 P þ j4 P þ j5 X þ j6 X
j1 þ j2 þ þ j6 ¼ 0
P ¼ Xþx
k ¼ XþxþM
After determining the form of the function S and the possible combinations of
the angles appearing in the expression of u, it is possible to identify those terms of
the disturbing function that bear significant contributions in the equations of

motion, and those terms that can be neglected [71].
Several methods have been proposed for computing the luni-solar perturbation.
An account of some of such methods is given below.
Roy [72] has studied the problem of the luni-solar perturbations affecting an
artificial Earth satellite in terms of changes induced in the orbital elements of the
satellite. Following Roy, the equations of motion of an artificial Earth satellite
subject to luni-solar perturbations are
@U @U @U
x00 ¼ y00 ¼ z00 ¼
@x @y @z
where x, y, and z are a set of three rectangular, non-rotating axes having their origin
in the centre O of the Earth, and
U ¼ VþF
is the force function. The force function, U, in turn, comprises the potential of
gravitation, V = GmE/r, due to a point-mass Earth, and the disturbing function, F,
which Roy expresses as follows

1 xxM þ yyM þ zzM 1 xxS þ yyS þ zzS
F ¼ G mM 3
þ mS 3
rM rEM rS rES
In the preceding expressions, G is the constant of gravitation, xM, yM, and zM are
the co-ordinates of the Moon, xS, yS, and zS are the co-ordinates of the Sun, x, y, and
z are the co-ordinates of the satellite, and mS, mM, and mE are the masses of,
respectively, the Sun, Moon, and Earth. In addition,
r 2 ¼ x2 þ y2 þ z 2
2
rES ¼ x2S þ y2S þ z2S
2
rEM ¼ x2M þ y2M þ z2M
2
rM ¼ ðxM xÞ2 þ ðyM yÞ2 þ ðzM zÞ2
rS2 ¼ ðxS xÞ2 þ ðyS yÞ2 þ ðzS zÞ2
Roy uses the disturbing function F given above in the Lagrange planetary
equations (see Sect. 3.2). Then, Roy computes approximately the perturbing effects
of the Moon and the Sun on an artificial Earth satellite, as follows

mM r 3 3r mS r 3 3r
M 2 1 S 2 1
m E rM 2rM m E rS 2rS
where r is less than rM or rS, and eM and eS are the magnitudes of the ratio of,
respectively, the lunar and the solar disturbing accelerations on the satellite to the
point-mass terrestrial acceleration on the satellite.
In addition, Roy computes eE, which is the magnitude of the ratio of the per-
turbing acceleration due to the second harmonic in the potential of gravitation of the
Earth to the point-mass terrestrial acceleration, as follows
r 2
E
E 3J 2
r
where J2 = –C20 = 0.001082636 (see Sect. 3.2) is the coefficient of the geopo-
tential expansion related to the Earth oblateness. This makes it possible to compute
eM, eS, and eE for artificial satellites orbiting the Earth at various altitudes. Roy finds
that the luni-solar perturbation is roughly equal in magnitude to the perturbation due
to the Earth oblateness for a satellite orbiting at an altitude of about 6.5 Earth radii
(41600 km), which is in line with the experimental results given by Blitzer [1].
Then, Roy sets
F ¼ FM þ FS
which means that the disturbing function (F) comprises a term (FM) due to the
Moon and another term (FS) due to the Sun. Roy considers two separate series
expansions for the two terms and truncates the two expansions to the first order.
This done, Roy uses separately the two expanded and truncated lunar and solar
terms of the disturbing function in the Lagrange planetary equations.
Sidi [73] expands the series of Legendre polynomials for 1/|r* − r| up to the
third term (k = 2), including the first term (k = 0). By so doing, the truncated
expansion of the disturbing function due to the presence of a third body becomes
r 2 1
l
F 1 þ 3 cos 2
c 1
r r 2
where ½(3 cos2c − 1) = P2(cos c), as has been shown in Sect. 3.2.
Kozai [74] computes the luni-solar perturbations upon the motion of an artificial
satellite revolving about the Earth in an orbit of high altitude, for example, upon the
motion of a synchronous satellite. Kozai, too, omits the first term (k = 0) from the
series of Legendre polynomials for 1/|r* − r| and assumes the constancy in time of
the orbital elements of the Sun and the Moon, except for the lunar ascending node
and perigee, whose motions are taken to be linear functions of time. By using
geocentric-equatorial co-ordinates x, y, and z, such that the x-axis is directed
towards the spring equinox, and the z-axis is directed towards the north pole, Kozai
expresses the three components of the position vector r of the satellite as a function
of the conventional orbital elements, as follows
x ¼ rðcos u cos X cos i sin u sin XÞ

y ¼ rðcos u sin X þ cos i sin u cos XÞ
z ¼ r sin u sin i
where u = x + / is the argument of latitude at epoch (see Chap. 1, Sect. 1.9).

According to the convention followed above, all quantities without an asterisk refer
to the artificial satellite, and those with an asterisk refer to the perturbing body.
As is easy to verify (see Chap. 1, Sect. 1.9), the expressions for x, y, and
z written above result from the matrix product R r, where the rotation matrix R and
the position vector r of the artificial satellite (in the geocentric perifocal reference
system) are respectively



and
2 3
r cos /
4 r sin / 5
0
The three components of the position vector r* of the perturbing body are
expressed by using its geocentric distance (r*), right ascension (a*) and declination
(d*), as follows
x ¼ r cos d cos a
y ¼ r cos d sin a
z ¼ r sin d
In the preceding expressions, r is the position vector of the artificial satellite and
r* is the position vector of the perturbing body, be it either the Moon or the Sun,
with respect to the Earth. Then, Kozai expresses the cosine of the angle, c, between
the position vectors r and r*, as follows
r r xx þ yy þ zz

cos c ¼ ¼ ¼ cos d cosðX a Þ cos u
rr rr
þ ½ cos d cos i sinðX a Þ þ sin d sin i sin u ¼ A cos u þ B sin u
where
A ¼ cos d cosðX a Þ
B ¼ cos d cos i sinðX a Þ þ sin d sin i
Then, in order to compute the Legendre polynomials P2(cos c), P3(cos c), …,
P5(cos c), which appear in the expression of the disturbing function F, Kozai
computes the second, third, fourth, and fifth power of cos c, as follows
cos2 c ¼ ðA cos u þ B sin uÞ2 ¼ A2 cos2 u þ B2 sin2 u þ 2AB sin u cos u

¼ 1=2 A2 þ B2 þ 1=2 A2 B2 cosð2uÞ þ AB sinð2uÞ

cos3 c ¼ 3=4ðA2 þ B2 ÞðA cos u þ B sin uÞ þ 1=4 A2 3B2 A cosð3uÞ

þ 1=4 3A2 B2 B sinð3uÞ
cos4 c ¼ 3=8ðA2 þ B2 Þ2 þ 1=2ðA2 þ B2 Þ½ðA2 B2 Þ cosð2uÞ þ 2AB sinð2uÞ

þ 1=8 A4 6A2 B2 þ B4 cosð4uÞ þ 1=2ðA2 B2 ÞAB sinð4uÞ
cos5 c 5=8ðA2 þ B2 ÞðA cos u þ B sin uÞ
In the expression of cos5c, all trigonometric functions whose arguments are 3u or

5u have been neglected. In Sect. 3.2, it has been shown that
P2 ðcos cÞ ¼ 1=2ð3 cos2 c 1Þ

P3 ðcos cÞ ¼ 1=2ð5 cos3 c 3 cos cÞ
P4 ðcos cÞ ¼ 1=8ð35 cos4 c 30 cos2 c þ 3Þ
P5 ðcos cÞ ¼ 1=8ð63 cos5 c 70 cos3 c þ 15 cos cÞ
Consequently, the four Legendre polynomials written above, from P2(cos c) to

P5(cos c), are
P2 ðcos cÞ ¼ 1=4½3ðA2 þ B2 Þ 2 þ 3=4½ðA2 B2 Þ cosð2uÞ þ 2AB sinð2uÞ

P3 ðcos cÞ ¼ 3=8 5 A2 þ B2 4 ðA cos u þ B sin uÞ þ 5=8½ðA2 3B2 ÞA cosð3uÞ
þ ð3A2 B2 ÞB sinð3uÞ
3 5 2
P4 ðcos cÞ ¼ ½35ðA2 þ B2 Þ2 40 A2 þ B2 þ 8 þ ½7 A þ B2 6
64 16
2 35 4
A B cosð2uÞ þ 2AB sinð2uÞ þ
2
½ðA 6A2 B2 þ B4 Þ
64
cosð4uÞ þ 4 A2 B2 AB sinð4uÞ
15
P5 ðcos cÞ ½21ðA2 þ B2 Þ2 28 A2 þ B2 þ 8ðA cos u þ B sin uÞ
64
In the expression of P5(cos5c), again, all trigonometric functions whose argu-

ments are 3u or 5u have been neglected.
By inserting the expressions of P2(cos c), …, P5(cos c) written above into

1 k
l X r
F¼ Pk ðcos cÞ
r k¼2 r
where F is the disturbing function, truncating the sum after k = 5, and replacing the
gravitational parameter l* = G(mc + m*) by n*2a*3, Kozai expresses the dis-
turbing function as follows

3 r r 2
2 2 a
F ¼ n br ½P2 ðcos cÞ þ P 3 ðcos cÞ þ P4 ðcos cÞ
r r r
r 3
þ P5 ðcos cÞ
r
where n* is the mean motion (introduced to the reader in Chap. 1, Sect. 1.3) of the
perturbing body, b = m*/(mc + m*), m* is the mass of the perturbing body, and mc
is the mass of the central body. In practice, the value of the factor b is equal to
unity, if the perturbing body is the Sun.
Kozai introduces the partial derivatives of disturbing function computed above
into the following Lagrangian planetary equations:

0 2 @F
a ¼
na @M

0 b2 @F b @F
e ¼
na e @M
4 na e @x
3

0 cos i @F 1 @F
i ¼
nab sin i @x nab sin i @X

0 cos i @F b @F
x ¼ þ
nab sin i @i na3 e @e

0 1 @F
X ¼
nab sin i @i

0 b2 @F 2 @F
M ¼n
na4 e @e na @a
where b = a(1 − e2)½. These equations must be integrated numerically, including the
terms related to other perturbations, for example, the Earth oblateness terms due to J2.
In an earlier article [75], Kozai has considered the luni-solar perturbations on the
orbit of an artificial satellite revolving about the Earth at low altitudes.
The solution of the Lagrangian planetary equations contains periodic, mixed,
and secular terms. Secular terms come from the part of the disturbing function
which does not contain periodic terms, and also from the part of the disturbing
function which does contain periodic terms, if the initial conditions are such as to
give rise to resonance phenomena. According to Gurfil [76], it is a common practice
in astrodynamics and celestial mechanics to compute secular and long-periodic
effects of a disturbing function on the orbit of an Earth satellite by the timescale
separation. This separation is obtained by averaging the Lagrangian planetary
equations. The averaging procedure provides expressions for the secular effects of
the first-order small perturbations on a satellite orbit, assuming that the variations of
orbital elements within a given time interval are of the second order.
The orbital elements computed on the basis of this assumption are called mean
orbital elements. The averaging, denoted by ⟨F⟩, of the disturbing function F is
performed as follows
Z2p
1
hF i ¼ FdM
2p
0
assuming that the disturbing function F is first-order small, that is,
F ¼ eF
with e 1, where ℱ is the non-dimensional potential. As a result of this

assumption, a timescale separation leaves the mean orbital elements unchanged, to
the first order, in the interval [0, 2p] of the true anomaly. The averaged differential
equations for the orbital elements of the Earth satellite result from substituting
a with ⟨a⟩, e with ⟨e⟩, etc., and also a′ with ⟨a⟩′, e′ with ⟨e⟩′, etc.
Domingos et al. [77] have applied the timescale separation technique by per-
forming a double average of the disturbing function over: (a) the mean motion of
the artificial satellite; and (b) the mean motion of the disturbing body. They expand
the Legendre polynomial series for the disturbing function up to the term corre-
sponding to k = 2 in the series. By so doing, and using below the same notation as
that used above, they express the disturbing function as follows

3
3
a 2 a r 21
F ¼ n 2 br 2
½ P2 ð cos c Þ ¼ n 2
ba
3 cos2 c 1
r r a 2

3 h i
1 a r 2
¼ n 2 ba2 3ðA cos u þ B sin uÞ2 1
2 r a
Then, Domingos et al. average the disturbing function over the eccentric anomaly
of the artificial satellite. To this end, they use the following identities (see Chap. 1,
Sect. 1.3):
1
ð1 e2 Þ2 sin Æ cos Æ e r
sin / ¼ cos / ¼ ¼ 1 e cos Æ
1 e cos Æ 1 e cos Æ a
M ¼ Æ e sin Æ dM ¼ ð1 e cos Æ ÞdÆ
where /, Æ, and M are, respectively, the true, eccentric, and mean anomaly, and
obtain the following expressions

r 2 2 1
cos / ¼ 1 þ 4e2
a 2

r 2 2 1
sin / ¼ 1 e2
a 2

r 2
sin / cos / ¼ 0
a

r 2 1
¼ 2 þ 3e2
a 2
Hence, the disturbing function, averaged over the mean anomaly of the artificial
satellite, becomes

3

3 a 2
hF i ¼ nbn 2 a2 A2 1 þ 4a2 þ B2 1 e2 þ e2
4 r 3
Then, Domingos et al. make the second average with respect to the mean
anomaly of the disturbing body. To do this, they consider the orbital elements of the
satellite constant during the process of average and obtain

a 1 3 2 15 4 2
A2 ¼ þ e þ e cos x þ cos2 i sin2 x
r 2 4 16

a 1 3 2 15 4 2
B2 ¼ þ e þ e sin x þ cos2 i cos2 x
r 2 4 16
Hence, Domingos et al. obtain the following expression for the double-averaged
disturbing function

hhF ii ¼ K 2 3 cos2 i 1 þ 3e2 3 cos2 i 1 þ 15e2 sin2 i cosð2xÞ
where

bn 2 a2 3 2 15 4
K¼ 1þ e þ e
16 2 8
Domingos et al. compute the partial derivatives of the double-averaged

disturbing function with respect to a, e, i, and x. By so doing, they obtain

@ hhF ii bn 2 a2 3 2 15 4
¼ 1þ e þ e ½2 3 cos2 i 1 þ 3e2 3 cos2 i 1
@a 8 2 8
þ 15e2 sin2 icosð2xÞ
@ hhF ii
¼ 6K e 3 cos2 i 1 þ 5e sin2 i cosð2xÞ
@e
@ hhF ii
¼ 3K 2 sinð2iÞ 3e2 sinð2iÞ þ 5e2 sin2 i cosð2xÞ
@i
@ hhF ii
¼ 30Ke2 sin2 i cosð2xÞ
@x
By introducing the preceding expressions into the Lagrangian planetary equa-

tions, Domingos et al. obtain
a0 ¼ 0
1

0 15bn 2 eð1 e2 Þ2 3 15 4
e ¼ 1 þ e 2 þ e sin2 i sinð2xÞ
8n 2 8

15bn 2 e2 3 15 4
i0 ¼ 1 1 þ e 2 þ e sinð2iÞsin x
16nð1 e2 Þ2 2 8

3bn 2 3 2 15 4
x0 ¼ 1 1þ e þ e ½5 cos2 i 1 þ e2 þ 5 1 e2 cos2 i
8nð1 e2 Þ2 2 8
cosð2xÞ

0 3bn 2 cos i 3 2 15 4 2
X ¼ 1 1þ e þ e 5e cosð2xÞ 3e2 2
8nð1 e2 Þ2 2 8

0 bn 2 3 2 15 4 2
M ¼ 1þ e þ e ½ 3e þ 7 3 cos2 i 1 þ 15 1 þ e2
8n 2 8
sin2 i cos2 x
The principal results found by Domingos et al. are shown below. Firstly, the
major semi-axis, a, of the artificial satellite does not change with time (a′ = 0). This
is because any gravitational field can be derived from a potential function, and
therefore mechanical energy is conserved [78]. On the other hand, in Chap. 1,
Sect. 1, it has been shown that the mechanical energy per unit mass of an artificial
satellite revolving about a main attracting body is E = –l/(2a), where l is the

gravitational parameter of the main body. Differentiating the preceding expression
with respect to time leads to
la0
E0 ¼
2a2
As the mechanical energy per unit mass of the satellite is constant in time (that
is, E′ = 0), so is its major semi-axis (a′ = 0). It is to be borne in mind that the result
a′ = 0 has been obtained in case of a doubly-averaged disturbing function.
Secondly, as has been shown by Costa et al. [79], there is a critical value, iC, of
the inclination, i, between the orbital plane of the artificial satellite and that of the
perturbing body, such that, if i > iC, then the eccentricity of the satellite orbit
increases. When this happens, a nearly circular orbit of a satellite about the Earth
becomes very elliptic. By contrast, when i < iC, then the satellite remains in its
nearly circular orbit. The case of a satellite in a nearly circular orbit is important,
because perturbations from other sources can change an exactly circular orbit into a
nearly circular orbit. Costa et al. have studied this problem using a doubly-averaged
disturbing function and a series expansion of Legendre polynomials carried out up
to P8(cos c). A short account of the results found by them is given below. With a
series expansion truncated after P2(cos c), the critical value, iC, of the orbital plane
of the satellite has been found to correspond to cos2i = 0.60, which yields
iC = 39°.2315. When the initial inclination, i0, of the orbital plane is sensibly below
the critical value (e.g. when i0 < 35°), then the eccentricity oscillates with a small
amplitude (less than 0.02 in most cases), which decreases rapidly with i0. In case of
more terms being used in the series expansion, some perturbations of shorter
periods appear, but the results are substantially the same as those found for an
expansion truncated after P2(cos c). If the initial inclination of the satellite orbit is
i0 = 0, then the orbital inclination remains constant, and the eccentricity has
oscillations of short period with very small amplitudes (about 0.05 for e0 = 0.5).
3.14 The Position of the Perturbing Body
The expressions considered in the preceding section give the Cartesian co-ordinates
(x*, y*, and z*) of the perturbing body with respect to a reference system having its
origin in the central body.
At first, let us take the Sun as the perturbing body. The method shown below,
due to Meeus [40], can be used when an accuracy of 0°.01 is sufficient. In this case,
a purely elliptical orbit of the Earth about the Sun can be assumed, so that the
perturbations of this orbit due to the Moon and the planets can be neglected.
Let JDE be the Julian Ephemeris Day of interest, computed as shown in
Sect. 3.4. Let T = (JDE−2451545.0)/36525 be the time, expressed in Julian
centuries, elapsed between the Julian ephemeris day of interest and J2000.0.
The quantity T must be computed with a sufficient number of decimals, in order to

avoid a loss of precision. Meeus points out that an error of 0.00001 in
T corresponds to an error of 0.37 day in time.
The geometric mean longitude of the Sun, with respect to the mean equinox of
date, is expressed by
L0 ¼ 280 :46645 þ 36000 :76983 T þ 0 :0003032 T 2
The mean anomaly of the Sun is
M ¼ 357 :52910 þ 35999 :05030 T 0 :0001559 T 2 0 :00000048 T 3
The eccentricity of the Earth orbit is
e ¼ 0:016708617 0:000042037 T 0:0000001236 T 2
Let C be the equation of centre (i.e. the true anomaly minus the mean anomaly)
of the Sun, which is computed by Meeus as follows
C ¼ þ ð1 :914600 0 :004817 T 0 :000014 T 2 Þ sin M

þ ð0 :019993 0 :000101 TÞ sinð2MÞ þ 0 :000290 sinð3MÞ
The true longitude (see Chap. 1, Sect. 1.9) of the Sun is kS = L0 + C, and the
true anomaly (see Chap. 1, Sect. 1.1) is / = M + C, by definition of the equation of
centre. The radius vector of the Sun, in astronomical units (AU), is
1:000001018ð1 e2 Þ
rS ¼
1 þ e cos /
If the value of rS is desired in km, the value computed above must be multiplied
by 1.495978707 108 (1 AU, [80]). In the fraction written above, the numerator
varies slowly with time. Meeus gives the following values
0.9997190 In the year 1800

0.9997204 1900
0.9997218 2000
0.9997232 2100
The longitude kS computed above is the true geometric longitude of the Sun with
respect to the mean equinox of date (which is not to be confused with the true
apparent longitude of the Sun). When it is necessary to compute the true geometric
longitude of the Sun with respect to the mean equinox of J2000.0, then it is possible
to use (with sufficient accuracy between the years 1900 and 2100) the following
expression
3.14 The Position of the Perturbing Body 423
ðkS ÞJ2000mod ¼ kS 0 :01397ðyear 2000Þ
The latitude uS of the Sun, with respect to the ecliptic of date, never exceeds the
value of 1″.2. For our purposes, uS can be set to zero.
Therefore, the right ascension and declination of the Sun can be obtained from
tan aS ¼ cos e0D sin kS = cos kS

sin dS ¼ sin e0D sin kS
As has been shown in Sect. 3.5, the mean obliquity of the ecliptic of date (e0D)
with respect to the mean equator of date is
e0D ¼ 8438100 :406 4600 :836769 T 000 :0001831 T 2 þ 000 :00200340 T 3

000 :000000576 T 4 000 :0000000434 T 5
(from Ref. [31]). The angle aS is always in the same quadrant as kS. If the numerator
and denominator on the right-hand side of the expression for aS are used in a
double-argument arctangent function (e.g., “atan2”), the proper quadrant will be
obtained. If aS is obtained in degrees, it can be converted to hours simply by
dividing by 15. Then aS is conventionally reduced to the range 0–24 h.
When the values of rS, kS, and e0D are known, the Cartesian co-ordinates of the
Sun (xS, yS, and zS) in the geocentric-ecliptic system xyz result from
xS ¼ rS cos uS cos kS rS cos kS

yS ¼ rS cos uS sin kS rS sin kS
zS ¼ 0
This approximation holds because, as indicated above, uS 0.

Hence, with reference to the following figure, the Cartesian geocentric-ecliptic
co-ordinates of the Sun can be converted to the geocentric Cartesian (true equator
and equinox of date) co-ordinates of the Sun as follows
2 3 2 3
XS rS cos kS
4 YS 5 ¼ R1 ð0D Þ4 rS sin kS 5
ZS 0
where, as has been shown in Sect. 3.4, R1(a) is the elementary rotation matrix
2 3
1 0 0
R 1 ð aÞ ¼ 4 0 cos a sin a 5
0 sin a cos a
The matrix multiplication indicated above yields
XS ¼ rS cos kS
YS ¼ rS sin kS cos e0D
ZS ¼ rS sin kS sin e0D
In the preceding expressions, the minus sign in front of e0D is due to the
clockwise rotation about the x-axis to transform the ecliptic xyz co-ordinates to the
equatorial XYZ co-ordinates, as shown in the preceding figure.
Another method for computing the Cartesian co-ordinates of the Sun in the
geocentric J2000-mean-equinox system (J2000mod), which has been shown in
Sect. 3.4, is suggested by Montenbruck and Gill [2]. This method is based on
appropriate mean orbital elements, which approximate the (apparent) elliptic orbit
of the Sun around the Earth for some decades around the year 2000. Such elements
are
a ¼ 149600000 km
e ¼ 0:016709
i ¼ 0 :0000
X þ x ¼ 282 :9400
M ¼ 357 :5256 þ 35999 :049 T
where T = (JD−2451545.0)/36525.0 is the number of Julian centuries elapsed from

J2000.0, and JD is the Julian Date of interest.
Taking the following value
e0J2000 ¼ 23 :43929111
for the obliquity of the ecliptic at J2000.0, and taking uS = 0, where uS is the
ecliptic latitude, the ecliptic longitude of the Sun (kS) and the Earth–Sun distance
(rS) result from
6892 72
kS ¼ X þ x þ M þ sin M þ sinð2M Þ
3600 3600
rS ¼ ½149:619 2:499 cos M 0:021 cosð2M Þ 106 km
and consequently the Cartesian co-ordinates of the Sun in the geocentric-equatorial

J2000-mean-equinox system (J2000mod) are
XS ¼ rS cos kS
YS ¼ rS sin kS cos e0J2000
ZS ¼ rS sin kS sin e0J2000
Once again, Montenbruck and Gill point out that the longitude kS, the latitude uS
and the radius vector rS of the Sun refer to the mean equinox and ecliptic of
J2000.0. In order to refer the same co-ordinates to the equinox of some epoch
T (measured in centuries from J2000.0), it is necessary to add a corrective term of
1 :3972 T
to the value of kS computed according to the Montenbruck–Gill method shown

above. There is no need to correct the value of the ecliptic latitude uS of the Sun,
since it varies by less than one arc minute within a full century [2].
Now, let us take the Moon as the perturbing body. To compute the position of
the Moon, Meeus [40] indicates a method having an accuracy of 10″ and 4″ in,
respectively, the longitude and latitude of the Moon. This method is based on the
five fundamental arguments (Delaunay variables) given in Sect. 3.5. However,
instead of using the longitude, X, of the ascending node, this method uses the mean
longitude, L0, of the Moon and obtains X from the difference
X ¼ L0 F
where F is the mean angular distance of the Moon from its ascending node. In the
same notation as that used so far, the fundamental arguments are expressed by
Meeus as follows
T3
l ¼ 134 :9634114 þ 477198 :8676313 T þ 0 :0089970 T 2 þ
69699
T4

14712000
T3
l0 ¼ 357 :5291092 þ 35999 :0502909 T 0 :0001536 T 2 þ
24490000
T3
E ¼ 93 :2720993 þ 483202 :0175273 T 0 :0034029 T 2
3526000
T4
þ
863310000
T3
D ¼ 297 :8502042 þ 445267 :1115168 T 0 :0016300 T 2 þ
545868
T4

113065000
T3
L0 ¼ 218 :3164591 þ 481267 :88134236 T 0 :0013268 T 2 þ
538841
T4

65194000
where, as shown in Sect. 3.4, T = (JDE−2451545.0)/36525, and JDE is the Julian

ephemeris day of interest. In addition, the following quantities are needed
A1 ¼ 119 :75 þ 131 :849 T

A2 ¼ 53 :09 þ 479264 :290 T
A3 ¼ 313 :45 þ 481266 :484 T
P P
It is necessary to compute the sums k (longitude) and r (distance) for the
Moon. To this end, the terms contained in
Pthe following table (from
P Ref. [40]) are used.
The argument (ai) of each sine (for k) and cosine (for r) function is a linear
combination of the four fundamental arguments D, l′, l, and F.
For example, the argument of the eighth line (except the headline) of the table is
2D þ ð1Þl0 þ ð1Þl þ 0F
P P
and the contributions of the eighth line to the to k and r are respectively
57066 sinð2D l0 lÞ
152138 cosð2D l0 lÞ
P P
D l′ l F k (coefficient of sin ai) r (coefficient of cos ai)
0 0 1 0 6288774 −20905355
2 0 −1 0 1274027 −3699111
2 0 0 0 658314 −2955968
0 0 2 0 213618 −569925
0 1 0 0 −185116 48888
0 0 0 2 −114332 −3149
2 0 −2 0 58793 246158
2 −1 −1 0 57006 −152138
2 0 1 0 53322 −170733
2 −1 0 0 45758 −204586
0 1 −1 0 −40923 −129620
1 0 0 0 −34720 108743
0 1 1 0 −30383 104755
2 0 0 −2 15327 10321
0 0 1 2 −12528
0 0 1 −2 10980 79661
4 0 −1 0 10675 −34782
0 0 3 0 10034 −23210
4 0 −2 0 8548 −21636
2 1 −1 0 −7888 24208
2 1 0 0 −6766 30824
1 0 −1 0 −5163 −8379
1 1 0 0 4987 −16675
2 −1 1 0 4036 −12831
2 0 2 0 3994 −10445
4 0 0 0 3861 −11650
2 0 −3 0 3665 14403
0 1 −2 0 −2689 −7003
2 0 −1 2 −2602
2 −1 −2 0 2390 10056
1 0 1 0 −2348 6322
2 −2 0 0 2236 −9884
0 1 2 0 −2120 5751
0 2 0 0 −2069
2 −2 −1 0 2048 −4950
2 0 1 −2 −1773 4130
2 0 0 2 −1595
4 −1 −1 0 1215 −3958
0 0 2 2 −1110
3 0 −1 0 −892 3258
2 1 1 0 −810 2616
(continued)
(continued)
P P
D l′ l F k (coefficient of sin ai) r (coefficient of cos ai)
4 −1 −2 0 759 −1897
0 2 −1 0 −713 −2117
2 2 −1 0 −700 2354
2 1 −2 0 691
2 −1 0 −2 596
4 0 1 0 549 −1423
0 0 4 0 537 −1117
4 −1 0 0 520 −1571
1 0 −2 0 −487 −1739
2 1 0 −2 −399
0 0 2 −2 −381 −4421
1 1 1 0 351
3 0 −2 0 −340
4 0 −3 0 330
2 −1 2 0 327
0 2 1 0 −323 1165
1 1 −1 0 299
2 0 3 0 294
2 0 −1 −2 8752
P
Now, it is necessary to compute the sum u (latitude) for the Moon. To this
end, the terms contained in the following table (from [40]) are used. Again, the
argument (ai) of each sine function is a linear combination of the four fundamental
arguments D, l′, l, and F.
P P
D l′ l F u (coeff. of sin D l′ l F u (coeff. of sin
ai) ai)
0 0 0 1 5128122 0 0 1 −3 777
0 0 1 1 280602 4 0 −2 1 671
0 0 1 −1 277693 2 0 0 −3 607
2 0 0 −1 173237 2 0 2 −1 596
2 0 −1 1 55413 2 −1 1 −1 491
2 0 −1 −1 46271 2 0 −2 1 −451
2 0 0 1 32573 0 0 3 −1 439
0 0 2 1 17198 2 0 2 1 422
2 0 1 −1 9266 2 0 −3 −1 421
0 0 2 −1 8822 2 1 −1 1 −366
2 −1 0 −1 8216 2 1 0 1 −351
2 0 −2 −1 4324 4 0 0 1 331
2 0 1 1 4200 2 −1 1 1 315
(continued)
(continued)
P P
D l′ l F u (coeff. of sin D l′ l F u (coeff. of sin
ai) ai)
2 1 0 −1 −3359 2 −2 0 −1 302
2 −1 −1 1 2463 0 0 1 3 −283
2 −1 0 1 2211 2 1 1 −1 −229
2 −1 −1 −1 2065 1 1 0 −1 223
0 1 −1 −1 −1870 1 1 0 1 223
4 0 −1 −1 1828 0 1 −2 −1 −220
0 1 0 1 −1794 2 1 −1 −1 −220
0 0 0 3 −1749 1 0 1 1 −185
0 1 −1 1 −1565 2 −1 −2 −1 181
1 0 0 1 −1491 0 1 2 1 −177
0 1 1 1 −1475 4 0 −2 −1 176
0 1 1 −1 −1410 4 −1 −1 −1 166
0 1 0 −1 −1344 1 0 1 −1 −164
1 0 0 −1 −1335 4 0 1 −1 132
0 0 3 1 1107 1 0 −1 −1 −119
4 0 0 −1 1021 4 −1 0 −1 115
4 0 −1 1 833 2 −2 0 1 107
P
The following expression is to be added to k:
þ 3958 sin A1 þ 1962 sinðk0 FÞ þ 318 sin A2

P
The following expression is to be added to u:
2235 sin k0 þ 382 sin A3 þ 175 sinðA1 F Þ þ 175 sinðA1 þ F Þ þ 127 sinðk0 lÞ
115 sinðk0 þ lÞ
The spherical co-ordinates of the Moon, with respect to the ecliptic, result from
P
k
kM ¼ L0 þ ðin degrees)
106
P
u
uM ¼ ðin degrees)
106
P
r
rM ¼ 385000:56 þ ðin km)
103
Dividing the sums by 106 or by 103 is necessary because the coefficients in the
tables shown above are given in units of 10−6 degree or 10−3 km.
Again, the mean obliquity of the ecliptic of date (e0D) with respect to the mean
equator of date is
e0D ¼ 8438100 :406 4600 :836769 T 000 :0001831 T 2 þ 000 :00200340 T 3

000 :000000576 T 4 000 :0000000434 T 5
(from Ref. [31]). Hence, the geocentric spherical ecliptic co-ordinates of the Moon
can be converted to geocentric Cartesian (true equator and equinox of date)
co-ordinates as follows
2 3 2 3
XM rM cos uM cos kM
4 YM 5 ¼ R1 ð0D Þ4 rM cos uM sinkM 5
ZM rM sin uM
where, as has been shown above, R1(a) is the elementary rotation matrix
2 3
1 0 0
R 1 ð aÞ 4 0 cos a sin a 5
0 sin a cos a
The matrix multiplication indicated above yields
XM ¼ rM cos uM cos kM
YM ¼ rM ðcos uM sin kM cos e0D sin uM sin e0D Þ
ZM ¼ rM ðcos uM sin kM sin e0D þ sin uM cos e0D Þ
This method for computing the position of the Moon implies the evaluation of
several terms contained in the tables given above. The evaluation, if done by hand,
requires much labour. In practice, the labour can be avoided by using the procedure
MOONPOS [81], written in the IDL scientific programming language and based on
the Meeus method. To the same end, Burnett has written an Excel® spreadsheet
[82], also based on the Meeus method.
Montenbruck and Gill [2] describe a method which has an accuracy of several
arc-minutes in the angles (uM and kM) and about 500 km in the lunar distance (rM).
The fundamental arguments (l, l′, F, D and L0) are computed by truncating the
expressions given above after the term containing the first power of T, as follows
l ¼ 134 :96292 þ 477198 :86753 T

l0 ¼ 357 :52543 þ 35999 :04944 T
F ¼ 93 :27283 þ 483202 :01873 T
D ¼ 297 :85027 þ 445267 :11135 T
L0 ¼ 218 :31617 þ 481267 :88088 T 1 :3972 T
where the term 1°.3972 T is subtracted from the linearised expression of L0, because
Montenbruck and Gill compute the co-ordinates of the Moon in the geocentric
J2000-mean-equinox system (J2000mod).
Then, Montenbruck and Gill compute the longitude of the Moon with respect to
the equinox and ecliptic of the year 2000, as follows
kM ¼ L0 þ 2264000 sinðlÞ þ 76900 sinð2lÞ 458600 sinðl 2DÞ þ 237000 sinð2DÞ

66800 sinðl0 Þ 41200 sinð2F Þ 21200 sinð2l 2DÞ 20600 sinðl þ l0 2DÞ
þ 19200 sinðl þ 2DÞ 16500 sinðl0 2DÞ þ 14800 sinðl l0 Þ 12500 sinðDÞ
11000 sinðl þ l0 Þ 5500 sinð2F 2DÞ
where the first two terms on the right-hand side describe the motion of the Moon in
an ellipse of eccentricity e = 0.055, and the remaining terms are due to various
perturbations. The lunar latitude results from
uM ¼ 1852000 sinðF þ kM L0 Þ þ 41200 sinð2FÞ þ 54100 sinðl0 Þ 52600 sinðF 2DÞ

þ 4400 sinðl þ F 2DÞ 3100 sinðl þ F 2DÞ 2500 sinð2l þ FÞ
2300 sinðl0 þ F 2DÞ þ 2100 sinðl þ FÞ þ 1100 sinðl0 þ F 2DÞ
where the first term on the right-hand side takes account of the inclination (about
5.1°) of the lunar orbit with respect to the ecliptic. The Earth–Moon distance
(measured in km between the two centres) is computed as follows
rM ¼ 385000 20905 cosðlÞ 3699 cosð2D lÞ 2956 cosð2DÞ 570 cosð2lÞ

þ 246 cosð2l 2DÞ 205 cosðl0 2DÞ 171 cosðl þ 2DÞ 152 cosðl þ l0 2DÞ
where terms smaller than 150 km have been neglected. Here, too, the spherical
ecliptic co-ordinates of the Moon (rM, uM and kM) can be converted to equatorial
Cartesian co-ordinates (XM, YM, and ZM), as has been shown above. It is to be noted
that the equatorial Cartesian co-ordinates computed by using the Montenbruck–Gill
method are the components of the position vector rM of the Moon in the geocentric
J2000-mean-equinox system.
Further methods than those shown above for computing the geocentric positions
of the Sun and Moon at a given time are given by Van Flandern and Pulkkinen [83],
Reda and Andreas [84], and Reda [85].
Online resources for the same purpose are provided by the United States Naval
Observatory [86], and Burnett [87]. The Naval Observatory also provides an
integrated package of subroutines and functions, called NOVAS, for computing
various commonly needed quantities in positional astronomy. This package is
available in FORTRAN, C, and Python. Details on this package can be found in
Ref. [88].
An example of application of the Meeus method for computing the Cartesian

co-ordinates of the Sun in the geocentric true equator and equinox of date reference
system is given below. Let 8 March 2012, noon UT1, be the date and time of
interest. By using Boulet’s method to compute the corresponding Julian Date (JD),
we set y = 2012, m = 3, d = 8, and h = 12. Then, the Julian Day Number (J0) and
the Julian Date (JD) result from
J0 ¼ 367y INTf1:75½y þ INTðm=12 þ 0:75Þg þ INTð275m=9Þ þ d þ 1721013:5

¼ 367 2012 INTf1:75 ½2012 þ INTð3=12 þ 0:75Þg þ INTð275 3=9Þ þ 8
þ 1721013:5 ¼ 2455994:5
h 12
JD ¼ J0 þ ¼ 2455994:5 þ ¼ 2455995:0
24 24
Hence, JDE results from
DT
JDE ¼ JD þ
86400
where DT is computed below by using the polynomial approximations indicated by

Espenak and Meeus [37], as follows
y ¼ year þ ðmonth 0:5Þ=12 ¼ 2012 þ ð3 0:5Þ=12 ¼ 2012:2083

s ¼ y 2000 ¼ 2012:2083 2000 ¼ 12:2083
DT ¼ 62:92 þ 0:32217s þ 0:005589s2 ¼ 67:686 s
Consequently, we have
67:686
JDE ¼ 2455995:0 þ ¼ 2455995:000783403
86400
Hence
JDE 2451545:0
T¼ ¼ 0:121834381476
36525
The geometric mean longitude of the Sun, with respect to the mean equinox of
date, is
L0 ¼ 280 :46645 þ 36000 :76983 T þ 0 :0003032 T 2 ¼ 4666 :59797939047

which is brought in the range [0°, 360°] by subtracting twelve times 360°:
4666 :59797939047 12 360 ¼ 346 :59797939047
The mean anomaly of the Sun is
M ¼ 357 :52910 þ 35999:05030 T 0 :0001559 T 2 0 :00000048 T 3

¼ 4743 :4511247009191
which is brought in the range [0°, 360°] by subtracting thirteen times 360°:
4743 :45111573436 13 360 ¼ 63 :4511247009191
The eccentricity of the Earth orbit is
e ¼ 0:016708617 0:000042037 T 0:0000001236 T 2 ¼ 0:0167034936
The equation of centre is
C ¼ þ ð1 :914600 0 :004817 T 0 :000014 T 2 Þ sin M

þ ð0 :019993 0 :000101 TÞ sinð2MÞ
þ 0 :000290 sinð3MÞ ¼ 1 :728112453543
The true longitude of the Sun is
kS ¼ L0 þ C ¼ 348 :326092358
The true anomaly is
/ ¼ M þ C ¼ 65 :1792371544621
The radius vector of the Sun, in astronomical units, is
1:000001018ð1 e2 Þ
rS ¼ ¼ 0:992760960276AU
1 þ e cos /
The mean obliquity of the ecliptic of date, with respect to the mean equator of
date, is
e0D ¼ 8438100 :406 4600 :836769 T 000 :0001831 T 2 þ 000 :00200340 T 3

000 :000000576 T 4 000 :0000000434 T 5 ¼ 23 :43769435337
The right ascension (aS) and declination (dS) of the Sun result from
Y ¼ cos e0D sin kS

X ¼ cos kS
a ¼ arctanðY=XÞ
If X\0 then aS ¼ a þ 180
If Y\0 and X [ 0 then aS ¼ a þ 360
Else aS ¼ a
dS ¼ arcsinðsin e0D sin kS Þ
In the present case, we find
Y ¼ cos e0D sin kS ¼ 0:185646791628

X ¼ cos kS ¼ 0:979315057919
a ¼ arctanðY=XÞ ¼ 10 :734075529
dS ¼ arcsinðsin e0D sin kS Þ ¼ 4 :616247987 ¼ 4 360 5800
Since Y < 0 and X > 0, then the right ascension of the Sun is
aS ¼ a þ 360 ¼ 349 :265924471 ¼ 23h 17m 04s
For comparison, the results which can be found in Ephemeris.com [89], where
the planet and Moon positions from NASA/JPL ephemeris files are used, are given
below:
aS ¼ 23h 17m 03s

dS ¼ 4 370 0700
Finally, the geocentric co-ordinates of the Sun (in the true equator and equinox
of date reference system) are
XS ¼ rS cos kS ¼ 0:972225757 AU
YS ¼ rS sin kS cos e0D ¼ 0:184302887 AU
ZS ¼ rS sin kS sin e0D ¼ 0:0798989791 AU
If the co-ordinates XS, YS, and ZS are desired in km, the values computed above
must be multiplied by 1.495978707 108, as has been shown above.
3.15 The Position of the Perturbing Body from NASA/JPL Ephemeris Files 435
3.15 The Position of the Perturbing Body from NASA/JPL

Ephemeris Files
The position of the perturbing body can be computed not only by using the methods
described in the preceding section but also by means of the NASA/JPL ephemeris
files. These files, named JPL DE(number) or simply DE(number), where DE stands
for Development Ephemeris, are models of the Solar System produced by the Jet
Propulsion Laboratory of Pasadena, California, for the purposes of space navigation
and astronomy. Such models provide estimates of the positions, velocities, and
accelerations of the principal celestial bodies of the Solar System, tabulated at equal
intervals of time, over a specified number of years. In particular, the tabulation
concerns the barycentric Cartesian co-ordinates of the Sun, the eight planets, Pluto,
and the geocentric co-ordinates of the Moon.
The data indicated above have been obtained by numerically integrating the
equations of motion for the celestial bodies, given a set of initial conditions coming
from high-precision observations. The integration was performed by solving
numerically the n-body problem. The initial conditions included both constant
quantities of the celestial bodies (e.g. their masses) and their initial state vectors, in
such a way as to best fit a large amount of observations. The best fit was obtained
by means of the least-squares technique.
Unless the desired ephemeris refers just to one of the tabulated times, an
interpolation is necessary. Therefore, the ephemeris data are given in the form of
numerical coefficients for Chebyshev polynomials. By evaluating these polyno-
mials, it is possible to determine the co-ordinates of the celestial bodies indicated
above with a precision of 0.01 arcsecond for the inner planets and of 0.1 arcsecond
for the outer planets. Of course, the actual precision of the results obtained depends
on the accuracy of the observations and the goodness of the least-squares fit.
Among the principal applications of the JPL ephemeris files are the navigation of
spacecraft through the Solar System for missions to the planets and ground-based
astronomy. The reason for using just Chebyshev polynomials to interpolate the
ephemeris data made available at the given times is shown below.
Following Mandel [90], let f(t) be a function of time defined on an interval [a,
b] of t. Let {t0, t1, …, tn} be a set of n + 1 distinct points (i.e. ti 6¼ tj, i 6¼ j), also
called nodes, contained in the interval [a, b]. The Lagrangian interpolation of f(t) is
defined as the unique polynomial p(t), of degree lower than or equal to n, which
satisfies the following condition
pð t i Þ ¼ f ð t i Þ
for i = 0, 1, …, n. If the (n + 1)th derivative of f(t) exists on (a, b) and if the nth
derivative of f(t) is continuous on [a, b], then the error of the Lagrangian interpo-
lation at the point t contained in [a, b] is
1
f ðtÞ pðtÞ ¼ qðtÞf ðn þ 1Þ ðsÞ
ðn þ 1Þ!
where s is contained in [a, b] and q(t) = (t − t0)(t − t1) … (t − tn-1)(t − tn).

In general, p(t) does not necessarily converge to f(t) as the number of nodes goes
to infinity. This is because the derivatives of f(t) may be large.
Let us consider the case of equidistant nodes t0, t1, …, tn. In this case, the
polynomial q(t) is small at the mid-point of the interval (t0, tn), but not small near to
the endpoints t0 and tn. The estimate of the error in the worst case (when t0 = a,
tn = b, tk − tk-1 = h, k = 1, 2, …, n) is
hh 1
j qð t Þ j 2h nh ¼ hn þ 1 n!
22 4
for t contained in [a, b]. A better interpolation than this can be achieved by making
a different choice of the nodes. In the hypothesis of |f(n+1)| M on the interval [a,
b], we want to choose the nodes t0, t1, …, tn in such a way as to make |f(t) − p(t)| as
small as possible. To this end, we have no control over f(t), because we do not know
this function. In addition, we cannot specify the point s in the interval [a, b]. What
we can do is to choose a set of nodes t0, t1, …, tn such that the maximum value of
|q(t)| should be minimum in the interval [a, b] of t. Thus, we take q(t) as a scaled
Chebyshev polynomial, as has been shown in Chap. 2, Sect. 2.10. This is because
the scaled Chebyshev polynomial Tn(x)/2n−1 is, of all polynomials of degree
n where the coeffient of xn is unity, the one having the smallest upper bound of its
magnitude within −1 x 1. The maximum magnitude of Tn(x) in this interval
being equal to 1, the upper bound mentioned above is 1/2n−1. Because of this
property, the scaled Chebyshev polynomials T(x)/2n−1 are those which deviate to
the least extent from zero. It follows that, if some error can be expressed as a
Chebyshev polynomial of degree n, then any other error expressible as a polyno-
mial of degree n and having the same coefficient of xn will have, within the interval
−1 x 1, a greater upper bound of its magnitude than the Chebyshev
expression of error.
To this end, as has been shown in Chap. 2, Sect. 2.10, we first operate a change
of variable, from a t b to −1 x 1, by defining the new variable x as
follows
2t ðb þ aÞ 2t ðtn þ t0 Þ
x¼ ¼
ba tn t0
and then take the nodes xk as follows

2k þ 1
xk ¼ cos p
2n þ 2
3.15 The Position of the Perturbing Body from NASA/JPL Ephemeris Files 437
for k = 0, 1, …, n. These nodes are just the zeros of the Chebyshev polynomial
Tn+1(x).
In other words, the Chebyshev interpolation on the interval −1 x 1 is
obtained as the Lagrangian interpolation of order n, where the nodes xk are chosen
as the zeros of Tn+1(x). Therefore, the function q(x) is chosen as follows
T n þ 1 ð xÞ
qð x Þ ¼
2n
so that |q(x)| 1/2n on the interval −1 x 1.

In summary, let f(t) be a function of t defined on the interval −1 t 1. As
has been shown above, in case of a definition interval [a, b] other than [−1, 1], it is
possible to change the variable, from t to x, such that −1 x 1.
The Chebyshev interpolant for f(t) on [−1, 1] is a polynomial p(t) of degree
n such as to match exactly f(t) at the zeros

2i þ 1
ti ¼ cos p
2n þ 2
(for i = 0, 1, …, n) of the (n + 1)th Chebyshev polynomial Tn+1(t).

The interpolant p(t) is a linear combination of Chebyshev polynomials such that
X
n
f ðtÞ pðtÞ ¼ ck T k ðtÞ
k¼0
whose coefficients ck result from
1 X n
c0 ¼ f ðt i Þ ðk ¼ 0Þ
n þ 1 i¼0
2 X n
ck ¼ T k ðti Þf ðti Þ ðk 6¼ 0Þ
n þ 1 i¼0
In the specific case of the JPL files, the planetary ephemerides are files of
Chebyshev polynomial coefficients fit to the barycentric Cartesian positions and
velocities of Mercury, Venus, the Earth–Moon barycentre, Mars, Jupiter, Saturn,
Uranus, Neptune, Pluto, and the Sun. The geocentric position of the Moon is
specified, so as to make it possible to determine, in combination with the
Earth-Moon barycentre position, the Solar System barycentric position of the Earth
or Moon. Following Fisher [47], each file of the JPL DE-series (e.g. DE 421,
described in Ref. [91]) begins with a header containing some lines of annotation,
the ephemeris time interval, the values of many constants assumed by the numerical
integration, and an index table which specifies the relative locations of the data for
the different bodies of the Solar System in the blocks of data which follow.
Each data block contains coefficients for Chebyshev polynomials, which specify the
value of each of the three Cartesian co-ordinates (X, Y, and Z) and, by numeric
differentiation, the value of each of the three velocity components for each celestial
body over the time interval of the data block, generally 32 days. All the dates and
coefficients are given in double precision. The positions of the celestial bodies are
integrated in astronomical units, but the polynomials are stored in units of kilo-
metres. The time-independent variable used in the integration is barycentric
dynamical time (TDB), which has been shown in Sect. 3.5, in fractional Julian
days.
Most JPL planetary ephemeris files include Chebyshev polynomial coefficients
fit to the lunar libration angles, which are integrated together with the planetary
positions. Many ephemeris files also fit to the 1980 IAU nutation series, which is
maintained in the files for backward compatibility.
In case of planetary positions being needed at a few specific times, the inter-
active website and telnet service of the JPL [92] can be used.
In case of planetary positions being needed at many times, the SPICE software
toolkit [93] can be used. This software toolkit is available in several programming
languages, including C, FORTRAN, and MATLAB, and for many platforms and
compilers.
3.16 The Radiation Pressure Due to the Sun
The solar radiation pressure is a perturbation due to the radiation emitted by the
Sun. This radiation exerts a pressure on the bodies hit by it, the amount of this
pressure depending on the properties of the surfaces of such bodies. This is because
the radiation incident on a body is partially absorbed, partially reflected specularly,
and partially reflected diffusively by the illuminated body.
Following Froideval [94], Max Planck found that the energy, E, of a quantum of
electromagnetic radiation is proportional to the frequency, f, of the radiation itself,
as follows
E ¼ hf
where h = 6.62606957 10−34 J s [95] is the Planck constant. In other words,

energy can be decomposed in discrete pieces, called quanta.
Einstein derived a relation between the energy, E, and the momentum, p, of a
quantum (photon) of light energy, as follows
E ¼ cp
where c = 299792458 m/s is the speed of light in vacuo.

3.16 The Radiation Pressure Due to the Sun 439
By combining together the two equations written above, there results
hf
p¼
c
The preceding expression, integrated over all the frequencies and further inte-
grated over the number of photons per unit area, yields the force per unit area. At
the distance from the Sun of one astronomical unit (AU), this force per unit area is
S0
F¼
c
where S0 is the total solar irradiance (TSI). Consequently, the total solar irradiance
is the amount of radiant energy emitted by the Sun over all the frequencies (or
wavelengths) which falls each second on the unit area (1 m2) exposed normally to
the Sun rays at the mean Sun–Earth distance in the absence of the Earth atmo-
sphere. Since the Earth orbit around the Sun is not circular but slightly elliptical,
then the intensity of the solar radiation received in a point outside the Earth
atmosphere varies with the square of the Earth–Sun distance.
On average, the total solar irradiance (also called solar constant) is Ssc = 1367
Watts per square metre (W/m2), according to Ref. [96]. This value varies by about
±3.4%, for the reason indicated above.
The Earth perihelion and aphelion occur around, respectively, 4 January and 5
July. Let Rav and R be, respectively, the mean Earth–Sun distance and the actual
Earth–Sun distance on the day of interest. The total solar irradiance, as a function of
R, is

2
Rav
S0 ¼ Ssc
R
According to Ref. [96], the quadratic factor in the preceding expression can be

2
Rav
¼ 1:00011 þ 0:034221 cos b þ 0:001280 sin b þ 0:000717 cosð2bÞ
R
þ 0:000077 sinð2bÞ
where b = 360°n/365, and n is the day of the year, starting at the 1st of January,
such that the 15th of January corresponds to n = 15, and the 15th of February
corresponds to n = 31 + 15 = 46. There are 366 or 365 days in a year depending
on whether the year is, or is not, a leap year. Stine and Harrigan [97] truncate the
expansion of (Rav/R)2 given above after the second term, as follows

360n
S0 ¼ Ssc 1 þ 0:034 cos
365:25
where the solar constant, Ssc, is taken equal to 1367.6 W/m2. The force per unit area
acting on a spacecraft placed at a distance d from the Sun is
Sd
F¼
c
where Sd, the total solar irradiance at a distance d from the Sun, is given by

2
D
Sd ¼ S0
d
and D is equal to 1 astronomical unit. It follows that the force per unit area is

S0 D 2
F¼
c d
Recently, Kopp and Lean [98] have found the value of 1360.8 ± 0.5 W/m2 for
the total solar irradiance during the 2008 solar minimum period. The last data
available at the moment of writing are given in the following figure, due to the
courtesy of the Laboratory for Atmospheric and Space Physics [99].
Several mathematical models have been proposed for computing the solar
radiation pressure acting on an artificial satellite.
The simplest of all them is the so-called cannon-ball model, which considers an
artificial satellite as a simple sphere. Following again Froideval [94], the acceler-
ation exerted by the solar radiation on such a type of satellite is
A
aSRP ¼ Pei ð1 þ vÞ uS
m
where P is the momentum flux of the solar radiation (see below), m is the reflectivity
coefficient of the satellite, A is the cross-sectional area of the satellite normal to the
Sun, m is the mass of the satellite, ei is the eclipse indicator (whose value is equal to
zero when the satellite in the full shadow of the Earth, and equal to unity when the
satellite is in full sunlight), and uS is the unit vector from the satellite to the Sun.
The minus sign on the right-hand side of the preceding equation is due to the fact
that the two vectors aSRP and uS are along the same straight line and oppositely
directed. This model is used in the Multi-Satellite Orbit Determination Program
(MSODP) of the University of Texas, where the momentum flux of the solar
radiation is expressed as follows

Ssc R 2
P¼
c r
where the solar constant, Ssc, is taken equal to 1367.2 W/m2, c is the speed of light
in vacuo (see above), R is one astronomical unit (whose value in km is given in
Sect. 3.14), and r is the distance of the satellite from the Sun. In the expression
written above, the value of the solar constant need not be re-computed as a function
of the Sun–satellite distance. This model is commonly used because of its
simplicity.
Other models, such as those described below, have been studied in order to
better approximate the satellites of complex shapes. In such cases, the satellite is
modelled as a combination of bodies comprising a box-shaped bus and a connected
solar array (see [100, 101]). This is because some satellites have large solar panels,
which make them unfit to be modelled as spheres. One of these satellites is illus-
trated in the following figure, due to the courtesy of Russian Space Systems [102].
This figure shows the GLONASS-M satellite, launched on the 10th of December of
2003, which has a surface of solar panels as large as 30.85 m2.
For satellites of this type, other models have been studied, one of which is the
Rockwell model, developed by Rockwell International, prime contractor of the GPS
satellites Block I and Block II.
The satellite-fixed reference system used in the Rockwell model is shown in the
following figure, where the image of the Earth is due to the courtesy of NASA.
Following Ziebart [103], the Cartesian, right-handed, satellite-fixed system of

co-ordinates X, Y, and Z has its origin in the nominal centre of mass of the satellite.
The antenna boresight is defined as the direction along the central bus longitudinal
axis, which points in the same direction as the navigation signal antennae. The
Z-axis coincides with the antenna boresight. The Y-axis is parallel to the solar panel
support boom arm, and the X-axis completes the right-handed set. Ziebart notes that
the satellite Z-axis points along the negative direction of the geocentric position
vector r of the satellite in an Earth-centred inertial (ECI) reference system, for
example, in the J2000.0 system. In other words, denoting by uZ the unit vector of
the satellite Z-axis, there results
r r
uZ ¼ ¼
jrj r
The satellite Y-axis is orthogonal to the vector p from the satellite to the Sun.
Let s be the position vector of the Sun in the same ECI reference system. Then,
there results
p¼sr
and consequently the satellite Y-axis is such that its unit vector uY is given by
uZ p
uY ¼
j uZ p j
The satellite X-axis is such that its unit vector uX forms, together with uY and uZ,
a Cartesian right-handed system. The positive direction of the satellite X-axis must
be in the same hemisphere as the Sun. This definition of the satellite-fixed
co-ordinate system specifies implicitly the satellite attitude [103].
According to this model, the satellite is represented by a combination of flat
surfaces and cylinders, and the solar rays are represented by a perfectly diffuse
beam and a perfectly specular beam. The force due to solar radiation is decomposed
into the following three components: a normal (FN) and a shear (FS) component due
to specular reflection, where the surface behaves like a perfect mirror; and a diffuse
(FD) component normal to the surface.
According to Fliegel et al. [104], in case of a flat surface (subscript f), these
components can be expressed as follows
AS0
FfN ¼ ð1 þ lmÞ cos2 h
c
AS0
FfS ¼ ð1 lmÞsin h cos h
c
2 AS0
FfD ¼ mð1 lÞcos h
3 c
In case of a cylindrical surface (subscript c), these components are

AS0 1
FcN ¼ 1 þ lm cos2 h
c 3
AS0
FcS ¼ ð1 lmÞsin h cos h
c
p AS0
FcD ¼ mð1 lÞcos h
6 c
where A is the area (in m2) of the cross section exposed to sunlight, S0 is the total
solar irradiance, c is the speed (in m/s) of light, m is the reflectivity coefficient, l is
the specularity coefficient, and h is the angle of incidence of the ray on the surface.
The preceding expressions are used in the standard models (S10 for Block I and
S20 for Block II).
The thermal models (T10 and T20) take account of the energy reradiated as heat
in both the satellite and the solar panels. By using the thermal models, the diffuse
component becomes for a flat surface
2 AS0
FfD ¼ mð1 lmÞcos h
3 c
and for a cylindrical surface
p AS0
FcD ¼ mð1 lmÞcos h
6 c
In the satellite-fixed XYZ reference system described above, the solar radiation
force F is decomposed into the three components FX, FY, and FZ.
These components can be expressed as a linear combination of trigonometric
functions of the angle e between the Z-axis and the unit vector up p/|p| pointing
towards the Sun, such that
cos e ¼ up uZ
According to Usai and Carpino [105], these expressions are
X
1 h i h i
ðX Þ ðX Þ ðX Þ ðX Þ
FX ¼ Ai sin Ki þ Ci cos Ki
i¼0
X
1 h i h i
ðZÞ ðZÞ ðZÞ ðZÞ
FZ ¼ Ai sin Ki þ Ci cos Ki
i¼0
In case of symmetrical satellites in nominal attitude (i.e. with their solar panels
orthogonal to the solar rays), there results FY = 0. However, a small FY term (called
Y-bias) can be produced by a misalignment of the solar panels or of the solar
sensors with respect to the panels or by the heat radiated by the louvres. The value
of the Y-bias is generally unknown and must be estimated together with a scale
factor for FX and FY.
The radial, transverse, and bi-normal components of the solar radiation accel-
eration vector a acting on the satellite are respectively
FZ
ar ¼
m
2 3
FZ 4 sin a cos k 5
a/ ¼
m 1 sin2 a sin2 k12
2 3
FZ 4 cos a 5
ah ¼
m 1 sin2 a sin2 k12
where m is the satellite mass, k is the satellite mean longitude, and a is the angle
between the unit vector up (from the satellite to the Sun) and the normal to the
orbital plane. By the way, as has been shown in Chap. 1, Sect. 1.10, the mean
longitude is k = M + X + x. The coefficients A(X) (X)
i , Ci , Ai
(Z),
and C(Z)
i , and the
(X) (Z)
multipliers Ki and Ki depend on the model used. In the complete T20 model, the
components FX and FZ of the solar radiation force are
FX ¼ 8:96 sin e þ 0:16 sinð3eÞ þ 0:10 sinð5eÞ 0:07 sinð7eÞ

FZ ¼ 8:43 cos e
where the coefficients are in units of 10−5 N [105].

The degree of goodness of the Rockwell model, for the purpose of describing the
real behaviour of the satellite in orbit, depends on the complexity of the satellite
modelled. With a mere box-wing model, some parts of the satellite, such as the
main antenna and the thrusters, are not taken into account. In addition, this model
takes into account only the first intersection of the rays of light with the surface of
the satellite and does not consider ageing phenomena that can change the specu-
larity properties of the surfaces [94].
These limitations have induced some authors, working at the Centre for Orbit
Determination in Europe (CODE), to develop the so-called Extended CODE Orbit
Model (ECOM). An outline of this model, which has been fully described by
Springer et al. [106], is given below. The ECOM is an empirical model. It expresses
the solar radiation acceleration in two different satellite-fixed reference systems.
The first is the XYZ system described above for the Rockwell model. The second is
the DYB system defined as follows. The D-axis is the direction satellite-Sun,
positive towards the Sun. The Y-axis is the same as that of the Rockwell model and
lies along the rotation axis of one of the solar panels. The B-axis completes the
right-handed system. The difference between the two reference systems (XYZ and
DYB) is shown in the following figure.
The IGS (International GNSS Service, where GNSS stands for Global
Navigation Satellite System) estimates parameters which best fit tracking data
received by GNSS ground stations forming a global network. By so doing, satellite
orbits have been computed with an accuracy of 2.5 cm for GPS satellites [107].
The ECOM has been developed on the basis of these observed data. With reference
to the preceding figure, let b0 be the elevation of the Sun above the orbital plane.
Let u0 and u be the arguments of latitude (u = x + /, see Chap. 1, Sect. 1.9)
relative to, respectively, the Sun and the satellite. Let uD, uY, uB, uX, and uZ be the
unit vectors relative to, respectively, the D-, Y-, B-, X-, and Z-axes shown above. In
accordance with Springer et al. [106], the acceleration vector (a) due to the solar
radiation can be written as follows
a ¼ aD uD þ aY uY þ aB uB þ aX uX þ aZ uZ
where the components of solar radiation acceleration are expressed as follows:
aD ¼ D0 þ DC2 cosð2b0 Þ þ DC4 cosð4b0 Þ

aY ¼ Y0 þ YC cosð2b0 Þ
aB ¼ B0 þ BC cosð2b0 Þ
aZ ¼ ½sinðuu0 Þ½Z0 þ ZC2 cosð2b0 Þ þ ZS2 sinð2b0 Þ þ ZC4 cosð4b0 Þ þ ZS4 sinð4b0 Þ
aX ¼ ½sinðuu0 Þ½X10 þ X1C cosð2b0 Þ þ X1S sinð2b0 Þ
þ ½sinð3uu0 Þ½X30 þ X3C cosð2b0 Þ þ X3S sinð2b0 Þ
The values of the three constants D0, Y0, and B0 were chosen to be
satellite-specific and are given below (first table). The values of the term Z0 were
chosen to be Block-type-dependent and are given below (second table). The values
of the constants DC2, DC4, YC, BC, ZC2, ZS2, ZC4, ZS4, X10, X1C, X1S, X30, X3C, and
X3S are given below (second table). All the values are taken from [106].
PRN Block D0 ( 10−9 m/s2) Y0 ( 10−9 m/s2) B0 ( 10−9 m/s2)

2 II −99.373 −0.6362 −0.0480
14 II −99.290 −0.9064 0.2510
15 II −98.985 −0.7084 0.4749
16 II −99.108 −0.6496 0.1170
17 II −99.010 −0.6604 0.0770
18 II −99.359 −0.8683 0.4783
19 II −99.850 −0.7057 0.1449
20 II −100.396 −0.6642 0.4997
21 II −99.477 −0.2592 −0.0996
1 IIa −91.088 −0.7458 0.4868
3 IIa −90.395 −0.5637 0.3960
4 IIa −90.502 −0.7856 0.2487
5 IIa −90.414 −0.7612 0.2309
6 IIa −90.354 −0.7589 0.3092
7 IIa −90.238 −1.0376 0.2241
8 IIa −93.342 −1.8394 0.7143
9 IIa −90.317 −0.7955 0.3569
10 IIa −89.546 −0.7819 0.1772
22 IIa −90.944 −0.7319 0.0179
23 IIa −78.592 −0.7440 1.0843
24 IIa −91.436 −1.0537 0.2214
25 IIa −90.785 −0.8556 0.3851
26 IIa −90.377 −0.9750 0.4144
27 IIa −90.291 −0.9482 0.4224
28 IIa −90.951 −0.8210 0.1303
29 IIa −91.015 −0.9078 0.5188
30 IIa −90.455 −0.8285 0.5409
31 IIa −90.370 −0.6269 0.6173
13 IIr −99.599 0.2801 1.6732
Parameters Estimate ( 10−9 m/s2) Formal error ( 10−11 m/s2)

DC2 −0.813 0.176
DC4 0.517 0.124
YC 0.067 0.104
BC −0.385 0.572
Z0 Block II 1.024 0.299
Z0 Block IIa 0.979 0.184
(continued)
(continued)
Parameters Estimate ( 10−9 m/s2) Formal error ( 10−11 m/s2)
ZC2 0.519 0.248
ZS2 0.125 0.149
ZC4 0.047 0.261
ZS4 −0.045 0.164
X10 −0.015 0.157
X1C −0.018 0.297
X1S −0.033 0.168
X30 0.004 1.655
X3C −0.046 3.118
X3S −0.398 1.773
According to Springer et al. [106], the performance of this model is better by

almost one order of magnitude than that of the Rockwell model. Their estimate of
the CODE model error is about 50 cm, which is to be compared with the about
300 cm in case of application of the Rockwell model.
The Jet Propulsion Laboratory (JPL) has developed two further empirical models
for solar radiation pressure acting on Earth satellites (Block IIa and Block IIr) used
for the Global Positioning System (GPS). The two models are fully described in,
respectively, [108, 109]. The latter model, which is an improvement of the former,
is described below. Bar-Sever and Kuang [109] base their model on four and
one-half years of precise GPS orbital data. They represent the satellite position by
means of the so-called EPS (e) angle, where EPS stands for Earth-Probe-Sun (or
Earth-Satellite-Sun). With reference to the preceding figure, e is the angle between
the Z-axis (pointing to the Earth) of the XYZ satellite-fixed system and the D-axis
(pointing to the Sun) of the DYB satellite-fixed system. For convenience of the
reader, the angle e is also shown in the following figure.
Bar-Sever and Kuang note that the XYZ system is that shown in the preceding
figure for Block II/IIa satellites, whereas, for Block IIr, the Z-axis is the same as that
shown above, but the sense of the X- and Y-axes is reversed (i.e. rotated by 180
degrees). Bar-Sever and Kuang start from the Fourier-series expansions (shown
above) of the T20 and T30 thermal models and, by means a procedure of trial and
error, arrive to the following expressions (in m/s2) of the three components of the
acceleration due to the solar radiation:
aX ¼ s105 ðAU=r Þ2 =m½S X1 sin e þ S X2 sinð2eÞ þ S X3 sinð3eÞ þ S X5 sinð5eÞ

þ S X7 sinð7eÞ
aY ¼ C Y0 þ 105 ðAU=r Þ2 =m½C Y1 cos e þ C Y2 cosð2eÞ
aZ ¼ s105 ðAU=r Þ2 =m½C Z1 cos e þ C Z3 cosð3eÞ þ C Z5 cosð5eÞ
where s is a dimensionless scaling factor (nominally unity), AU is the astronomical

unit (whose value in km is given in Sect. 3.14), r (in km) is the distance between
the satellite and the Sun, and m (in kg) is the mass of the satellite. The dimension of
the Fourier expansions enclosed between square brackets is 10−5 newtons. Since the
solar radiation pressure acting on a satellite depends on time variations affecting the
solar flux and the satellite mass, then Bar-Sever and Kuang include these small
un-modelled variations into the parameter estimation in each individual satellite arc
fit through an overall scale factor [109]. The coefficients of the Fourier series are
not, all of them, constant for a given satellite and model. Some of these coefficients,
as is the case with C_Y1 and S_X2, depend on the angle (b0) between the Sun–Earth
line and the orbital plane of the satellite. For the two coefficients named above,
Bar-Sever and Kuang found a dependency on b0 of the following type
F ðb0 Þ ¼ A þ B sin b0 þ C= sin b0 þ D cos b0
The parameters A, B, C, and D were estimated in a least-squares fit, and the

results are given below in two separate tables, for the GSPM.04a model and for the
GSPM.04b model.
GSPM.04a model—value of parameters in 10−5 Newtons
Parameter Block IIa (10−5 N) Block IIr (10−5 N)

S_X1 −8.982 10.931
S_X2 −0.0219 0.1279
S_X3 0.0151 0.2767
S_X5 0.1040 −0.2045
S_X7 0.0038 0.0568
C_Z1 −8.6044 −11.6408
C_Z3 0.0158 0.0627
C_Z5 0.0553 0.0674
C_Y2 0.01729 −0.0067
C_Y1 0.0091 + 0.0539 sin b0 + 0.0265/sin b0 0.0010 − 0.0199 sin b0 − 0.0107/sin b0
GSPM.04b model—value of parameters in 10−5 Newtons
Parameter Block IIa (10−5 N) Block IIr (10−5 N)

S_X1 −8.9820 10.9310
S_X2 −0.0509 + 0.0002 sin b0 + 0.0002/ 0.0172 + 0.0022 sin b0 − 0.0016/
sin b0 + 0.0407 cos b0 sin b0 + 0.1477 cos b0
S_X3 0.0045 0.2476
S_X5 0.1060 −0.2283
S_X7 0.0028 −0.0140
C_Z1 −8.6044 −11.6411
C_Z3 0.0225 0.0583
C_Z5 0.0543 0.0571
C_Y2 0.0175 −0.0064
C_Y1 0.0271 + 0.0459 sin b0 + 0.0302/ −0.0195 − 0.0172 sin b0
sin b0 −0.0252 cos b0 −0.0119/sin b0 + 0.0272 cos b0
The following remarks apply to both of the tables shown above. First, the
satellite-fixed XYZ system is that which is illustrated in the preceding figure for
Block II/IIa satellites, whereas, for Block IIr, the Z-axis is the same as that shown
above, but the sense of the X- and Y-axes is reversed (i.e. rotated by 180°). Second,
for |b0| < 0.25307 radians (14°.5), b0 is to be fixed to 0.25307 when b0 is positive,
and to −0.25307 when b0 is negative.
Finally, the so-called UCL model is described below. This model is due to
Ziebart [103] and has been applied by Munyamba [110] to the complex shape of the
GLONASS-M satellite, in such a way as to improve the Rockwell box-wing model
described above. The analytical modelling used for solar radiation pressure is
illustrated in the following figure.
Let A, h, l and m be, respectively, the area of the flat surface of the satellite, the
angle of incidence of the rays, the specularity coefficient and the reflectivity
coefficient of the surface. Following Ziebart [103], the area of the surface orthog-
onal to the direction of the Sun rays is A cos h.
Let F = S0/c be the force per unit area acting at the distance from the Sun of one
astronomical unit (AU), where S0 is the total solar irradiance (TSI) and c is the
speed of light in vacuo, as has been shown at the beginning of this section.
Firstly, the magnitude of the force vector FA due to absorbed light is
S0
FA ¼ A cos h
c
The shear (subscript S) and normal (subscript N) components of the force vector
FA are respectively
S0
FAS ¼ A cos h sin h
c
S0
FAN ¼ A cos2 h
c
Secondly, let us consider the reaction force vector FSR due to the specularly
reflected light. Of the total amount (S0/c) A cos h of the reaction force, the reflected
part is m (S0/c) A cos h, where m is the reflectivity coefficient of the material. The
reflected part, multiplied by the specularity coefficient l of the material, yields the
part of specularly reflected light (lm (S0/c) A cos h), whose reflection angle, with
respect to the normal to the surface, is equal to the incidence angle. The shear
(subscript S) and normal (subscript N) components of the reaction force vector FSR,
due to the specularly reflected light, are respectively
S0
FSRS ¼ lm A cos h sin h
c
S0
FSRN ¼ lm A cos2 h
c
Thirdly, the following figure shows the part of light which is reflected diffusely.
The amount of light which is reflected diffusely is
S0
mð 1 l Þ A cos a
c
Therefore, the reactive force vector FDF directed along the normal, due to the
light reflected diffusely, results from integrating over a hemisphere.
This yields the magnitude of FDF, that is,
2 S0
FDF ¼ mð1 lÞ A cos h
3 c
Taking together the shear (subscript S) and normal (subscript N) components of

the total force vector F leads to
S0
FS ¼ ð1 lmÞ A cos h sin h
c

S0 2
FN ¼ A ð1 þ lmÞcos h þ mð1 lÞ cos h
c 3
These expressions hold in case of a flat plate hit by the rays of the Sun. In case of
a cylindrical surface, the corresponding expressions are those given below
S0
FS ¼ ð1 lmÞ
A cos h sin h
c

S0 1 p
FN ¼ A 1 þ lm cos h þ mð1 lÞ cos h
c 3 6
where A is the cross-sectional area presented to an observer whose line of sight is

perpendicular to the axis of the cylinder.
Ziebart’s model [103] considers the radiation emitted by the Sun as a plane
wave-front, constituting a pixel array, at the Sun–spacecraft distance.
As shown in the preceding figure, each pixel of this array represents a ray of light,
whose sectional area is equal to the area of the pixel in the array plane. These pixels
are projected from the plane wave-front towards the spacecraft along a direction
which is parallel to the Sun–spacecraft vector. The pixel array is used to sample a
computer model of the spacecraft surface to determine, as a function of the EPS (e)
angle, which of all the parts of the spacecraft is illuminated and which is in shadow.
The solar panels are oriented perpendicularly to the line joining the satellite with
the Sun, whereas the antenna boresight is oriented towards the centre of the Earth,
not towards the Sun, in such a way as to form the EPS (e) angle with the Sun–
satellite line.
Ziebart also notes that, generally speaking, the following elements are needed in
the development of analytical models of solar radiation pressure:
• a structural description of the spacecraft, including its mass, dimensions, and
optical properties of the surface materials;
• the attitude of the spacecraft in the course of its trajectory; and
• the parameters which describe the environment in which the spacecraft operates
during its lifetime.
Georgevic [111] has given the values of the reflectivity coefficients l and m with
reference to the materials used in the Mariner 9 Mars orbiter spacecraft. Some of
such values are given in the following table (from [111]).
Spacecraft component l m
Solar panels 1.00 0.21
Propulsion module 0.67 0.64
High-gain antenna 0.67 0.30
Low-gain antenna 0.67 0.70
Solar panel outriggers 1.00 0.80
Solar panel end beams 1.00 0.14
The components of the solar radiation force computed above are relative to a
non-inertial body-fixed reference system XYZ (or DYB), having its origin in the
barycentre of the spacecraft. By contrast, the equations of motion are written in an
inertial reference system (e.g. in the heliocentric ICRS described in Sect. 3.8). In
order to transform such components to those relative to the heliocentric ICRS, it is
necessary to compute the angular velocity vector x.
The details can be found, for example, in Georgevic [112].
3.17 The Eclipse Factor
As has been shown in the preceding section, a satellite is subject to a solar radiation
force only when it is illuminated by the Sun.
When a satellite passes through the shadow of the Earth (called umbra in the
following figure), there is no solar radiation force acting directly upon it. Therefore,
it is necessary to compute the so-called eclipse factor, that is, the fraction of orbital
period in which a satellite is in the shadow of the Earth.
Let T and t2 − t1 (where 0 t2 − t1 < ½T) be, respectively, the revolution

period and the eclipse time of an Earth satellite. The eclipse factor, ef, is defined as
follows
t2 t1
ef ¼
T
The calculation of the eclipse factor is also important in the thermal analysis of
the satellite. This is because this analysis is carried out by evaluating the effects
produced by the various causes that produce heating. One of such causes is
certainly the amount of solar radiation incident directly on the satellite, which is
one of the main sources of energy. The evaluation of the solar radiation requires,
in turn, a previous determination of the eclipse times. In most practical cases, the
3.17 The Eclipse Factor 455
eclipse times need not be evaluated through the complex analytical formulation
that follows from taking into account all the perturbations to the central force field.
It is, on the contrary, expedient to use simple methods to calculate quickly the time
per orbital period during which the satellite passes through the shadow of the
Earth.
The present section concerns such methods. They use simple equations and lead
to solutions approximate enough for the purposes of a definition study. The
methods described below are based, in accordance with Dreher [113], on the fol-
lowing assumptions:
(a) The Earth and the Sun are perfect spheres of diameters DE and DS, respectively.
(b) All perturbations to Keplerian orbits are negligible.
(c) The Earth–Sun distance, rS, is constant.
(d) The relative motion between the Sun and the Earth during the time of satellite
transit through the shadow of the Earth is negligible.
Cunningham [114, 115] adds the following ones:
(e) For the calculation of eclipse times, t2 − t1, the Earth is considered as a sphere
with a radius of 6371 km. This means that the geometrical flattening of the
Earth and consequently the apparent contour of the Earth are not taken into
consideration. Refraction effects on the terrestrial atmosphere are also
neglected.
(f) The shadow of the Earth is approximated by a circular cylinder. Such an
approximation is valid because the umbra (i.e. the conical total shadow pro-
jected from the Earth on the opposite side of the Sun, where the intensity of the
solar radiation is zero) is one terrestrial diameter in width and approximately
106 terrestrial diameters in length, as shown in the following figure (courtesy of
Wikimedia [116]), and then extends well beyond the mean Earth–Moon dis-
tance, which is about thirty Earth diameters. Therefore, the umbra width
changes very little in the vicinity of the Earth.
According to Robertson [117], for satellite orbits within an altitude of one Earth
radius, this approximation introduces an error in the calculation of eclipse times of
the order of magnitude of 1% of the satellite orbital period. This approximation is
also good for Earth satellites in geosynchronous orbits, the shadow times resulting
from it being slightly longer than the actual shadow times. Consequently, no
attempt is made at first to consider the penumbra (i.e. the partial-shadow region
placed between the umbra and the full-light region) where the satellite is illumi-
nated by a portion of the solar disc.
Under the hypotheses specified above, two methods, which are due to, respec-
tively, Cunningham [114, 115] and Escobal [9, 118], are described below, by which
first-approximation values of eclipse times for either circular or elliptic orbits of
Earth satellites can be computed.
Should the case require more accuracy than is possible with these methods, the
eclipse time analysis may be divided into two sections. In the first, the true
anomalies relating to the points of, respectively, shadow entrance and shadow exit
are determined by means of the simple mathematical model based on the
assumptions mentioned above. In the second, the values of the true anomalies found
in the first section are used as an initial approximation to be refined by including the
effects produced by the perturbations, as the sequel will show.
Since the Earth shadow is assumed to be cylindrical, the intersection of the orbit
plane with the shadow forms a semi-ellipse (shown in the following figure), which
can be constructed if the angle between the Sun–Earth direction and the orbit plane
is known.
This angle can easily be computed if the Sun–Earth unit vector (s) and the
satellite position (r) and velocity (v) vectors are known at some time t, as indicated
below. Following Robertson [117], the unit vector uz perpendicular to the orbital
plane is given by
rv
uz ¼
j r vj
The angle (d) between the orbital plane and the Sun–Earth unit vector (s) results
from
p
cos d ¼ uz s
2
that is,
p
d¼ arccosðuz sÞ
2
The preceding figure also shows that the major semi-axis (aU) and the minor
semi-axis (bU) of the shadow semi-ellipse are, respectively, aU = RE/sin d and
bU = RE, where the subscript U relates to the shadow (umbra) of the Earth.
The following figure shows that, in the particular case of a satellite in a circular
orbit of radius r around the Earth, the eclipse time is the time spent by the satellite
to pass through an angle of 2h, where h is computed by considering the intersec-
tions of the shadow semi-ellipse with the circular orbit of the satellite.
Therefore, h has the following expression

"
1 # " 1 #
RE a2U r 2 2 R2E r 2 sin2 d 2
h ¼ arcsin ¼ arcsin
r a2U R2E r cos d
The time t2 − t1, during which the satellite passes through the shadow of the
Earth, is then

3 12
2h r
t2 t1 ¼ T ¼ 2h
2p lE
where T = 2p (r3/lE)½ is the orbital period of the satellite and lE = GME is the
gravitational parameter of the Earth. Then, the eclipse factor, in case of circular
orbits, is ef = h/p.
In case of satellites revolving around the Earth in elliptical orbits, two methods
for computing the eclipse times (or the eclipse factors) are described below.
According to Cunningham’s method, as the Earth has been assumed to be
perfectly spherical and all perturbing effects have been neglected, the satellite orbit
is an ellipse with the centre of the Earth at a focus:
að 1 e 2 Þ
r¼
1 þ e cosðw cÞ
where r is the radius vector, a is the major semi-axis, e is the eccentricity, w is the
true anomaly measured from the projection of the unit Sun–Earth vector onto the
orbital plane, and c is the angle formed by the line of apsides and the projection of
the unit Sun–Earth vector onto the orbital plane, as shown in the following figure,
where, for simplicity of representation and without loss of generality, the Sun–Earth
direction is assumed to lie in the orbital plane (d = 0).
In this case, the shadow semi-ellipse degenerates into a rectangular strip, because
the intersection between the umbra region (assumed to be cylindrical) and the
orbital plane is just a rectangle. The positive direction of w is determined by the
manner in which c is defined. In this case, c is defined with respect to the line of
apsides at the point of perigee such that 0 c p, as shown in the following
figure. There is then complete symmetry about the line of apsides.
In this manner, the orientation of the elliptic orbit is exactly defined by specifying
the angles c and d. Of course, when d = p/2 the projection of the Sun–Earth unit
vector onto the orbital plane yields a point, not a line, which leaves c undefined.
For each orbit, there is an upper limit d′ of d (where d′ < p/2) such that for
d
d′ the satellite is in sunlight during the whole orbital period. It is evident that
any value of d > 0 can be obtained by simply rotating the orbital plane about the
line passing through the centre of the Earth perpendicular to the Sun–Earth unit
vector.
Precisely, d′ is defined as that value of d which occurs when the orbit is tangent
to the shadow of the Earth. The preceding figure shows that when 0 d < d′ the
satellite spends a portion t2 − t1 (such that 0 t2 − t1 < T/2) of its orbital period
T in shadow; likewise, when d′ d p/2 the satellite is in sunlight for the entire
orbit.
The projection of r (in this case, the radius vector lying along the continuation of
the projection of the Sun–Earth unit vector onto the orbital plane) onto the plane
normal to the Sun–Earth unit vector is equal to the Earth radius when d = d′, as
shown in the preceding figure. Hence we have

0 að 1 e 2 Þ
r sin d ¼ sin d0 ¼ RE
1 þ e cosðw cÞ
that is,
að1 e2 Þsin d0
¼ RE
1 þ e cosðp cÞ
from which d′ can be determined.

The final step is the determination of the true anomalies w1 (of shadow entrance)
and w2 (of shadow exit) as functions of d in the range 0 d d′. This may easily
be done by computing the intersections of the shadow semi-ellipse with the satellite
orbit, as shown in the following figure.
The computation may be done by writing the equation of the shadow

semi-ellipse in polar co-ordinates (r, w) with origin in the occupied focus (i.e. in the
centre of the Earth) of the elliptic orbit of the satellite:
RE
r¼ 12
1 cos2 w cos2 d
and the equation of the satellite orbit:
að 1 e 2 Þ
r¼
1 þ e cosðw cÞ
By intersecting the shadow semi-ellipse with the satellite orbit, we find
a 1
f ðwÞ 1 þ e cosðw cÞ 1 e2 1 cos2 w cos2 d 2
RE
The zeros of the shadow function f(w) may be determined through a suitable
numerical technique, as will be shown below. From geometrical considerations on
the eclipse phenomenon (see the preceding figure), it is known that one (i.e. w1) of
the two physically meaningful zeros lies in the second quadrant, and the other one
(i.e. w2) lies in the third quadrant, because this is the sole domain where an eclipse
can take place. The search for w1, therefore, is to be done in the interval going from
p/2 to p; and the search for w2, from p to 3p/2. Once the values of w1 and w2 are
known, the time t2 − t1, spent by the satellite to pass through the shadow of the
Earth, can be computed in more than one way. Cunningham [114, 115] proposes
the following expression
T h 1 i
t2 t1 ¼ Tef ¼ e 1 e 2 2 ð p1 p2 Þ þ 2ð q2 q1 Þ
2p
where n = 1 or 2, T = 2p(a3/lE)½ is the orbital period of the satellite, and

"
1
#
sinðwn cÞ 1e 2 wn c
pn ¼ qn ¼ arctan tan
1 þ e cosðwn cÞ 1þe 2
Alternatively, the shadow time t2 − t1 can be computed as suggested by Escobal

[9, 118], by means of the eccentric anomalies Æ1 and Æ2 relating to the points of,
respectively, shadow entrance and shadow exit, as will be shown below. Once the
values of w1 and w2 are known, the values of the corresponding true anomalies /1 and
/2 (measured from the perigee direction) result from /1 = w1 − c and /2 = w2 − c.
Hence, remembering the expressions of Chap. 1, Sect. 1.3, the values of the corre-
sponding eccentric anomalies Æ1 and Æ2 result from
1
ð1 e2 Þ2 sin / e þ cos /
sin Æ ¼ cos Æ ¼
and the shadow time results from

12

a3 sin Æ 2 sin Æ 1
t2 t1 ¼ arctan arctan eðsin Æ 2 sin Æ 1 Þ
lE cos Æ 2 cos Æ 1
where the correct quadrant of Æ is to be determined by examination of the

numerator and denominator of the arctangents, and the difference of the arctangents
should be a positive number [118].
The shadow time t2 − t1 can also be computed by first evaluating the radii r1 and
r2, joining the centre of the force field (i.e. the geometric centre of the Earth) with
the points of shadow entrance and shadow exit, respectively, and the chord
c joining these points, as indicated below:
að 1 e 2 Þ
r1 ¼
1 þ e cos /1
að 1 e 2 Þ
r2 ¼
1 þ e cos /2
1
c ¼ r12 þ r22 2r1 r2 cosð/2 /1 Þ 2
Then, by applying Lambert’s theorem (see Chap. 1, Sect. 1.5), the shadow time
is

12
T a3
t2 t1 ¼ ½ða sin aÞ ðb sin bÞ ¼ ½ða sin aÞ ðb sin bÞ
2p lE
where T = 2p(a3/lE)½ is the orbital period of the satellite, and a and b are,
respectively, Lambert’s first and second angle, expressed by
a r þ r þ c12

r1 þ r2 c2
1
1 2 b
sin ¼ sin ¼
2 4a 2 4a
The angles d and c, for the day of interest, can be computed as follows. First, we
determine the position vector of the Sun in the geocentric ICRS of epoch, denoted
by XYZ. Let XS, YS, and ZS be the three components of the position vector of the Sun
with respect to XYZ.
Then, we determine the components of the Sun–Earth unit vector s with respect
to XYZ. Let sX, sY, and sZ be these components. They are given by
XS YS ZS

sX ¼ 12 sY ¼ 12 sZ ¼ 12
XS2 þ YS2 þ ZS2 XS2 þ YS2 þ ZS2 XS2 þ YS2 þ ZS2
Then, we transform sX, sY, and sZ to sx, sy, and sz, where sx, sy, and sz are the three
components of the Sun–Earth unit vector s with respect to the geocentric perifocal
reference system xyz.
As has been shown in Chap. 1, Sect. 1.9, xyz is a reference system which has its
origin in the centre of the Earth; the fundamental plane xy is the plane to which the
motion of the satellite is confined. Let ux, uy, and uz be the three unit vectors along
the x, y, and z axes of the perifocal reference system. Of these unit vectors, ux and
uy are along the two Cartesian axes x and y which are contained in the orbital plane,
oriented so that ux points towards the perigee, and uy is 90° ahead of ux in the
direction of orbital motion; whereas uz is the unit vector along z, perpendicular to
the orbital plane, oriented so that ux uy uz should be a right-handed reference
system.
To this end, uz points towards the same direction as the moment of momentum
per unit mass of the satellite h = r r′.
As has been shown in Chap. 1, Sect. 1.9, the transformation from sX, sY, and sZ
to sx, sy, and sz is performed by means of the following rotation matrix

cos X cos x sin X sin x cos i sin X cos x þ cos X sin x cos i sin x sin i

cos X sin x sin X cos x cos i sin X sin x þ cos X cos x cos i cos x sin i

sin X sin i cos X sin i cos i
where X, x, and i are, respectively, the right ascension of the ascending node, the
argument of perigee, and the inclination of the satellite orbital plane with respect to
the equatorial plane. By multiplying this rotation matrix by the column vector
s sX uX þ sY uY þ sZ uZ
according to the rules of matrix multiplication, there results
sx ¼ ðcos X cos x sin X sin x cos iÞsX þ ðsin X cos x þ cos X sin x cos iÞsY
þ ðsin x sin iÞsZ
sy ¼ ð cos X sin x sin X cos x cos iÞsX þ ð sin X sin x þ cos X cos x cos iÞsY
þ ðcos x sin iÞsZ
sz ¼ ðsin X sin iÞsX þ ð cos X sin iÞsY þ ðcos iÞsZ
In the event of the values of sx and sy being, both of them, equal to zero, d = 90°
and therefore the satellite is constantly exposed to the Sun (t2 − t1 = 0). Otherwise
2 !12 3
s2
d ¼ arctan4 2 z 2 5
sx þ sy
and c depends on the values of sx and sy, as follows
sx \0 and sy ¼ 0 c ¼ 0
sx \0 and sy [ 0 c ¼ arctanðsy =sx Þ
sx ¼ 0 and sy [ 0 c ¼ 90
sx [ 0 and sy [ 0 c ¼ 180 arctanðsy =sx Þ
sx [ 0 and sy ¼ 0 c ¼ 180
sx [ 0 and sy \0 c ¼ 180 arctanðsy =sx Þ
sx ¼ 0 and sy \0 c ¼ 90
sx \0 and sy \0 c ¼ arctanðsy =sx Þ
In addition to the method described above, Escobal [9, 118] and, more recently,
Neta and Vallado [119] have proposed to determine the true anomalies /1 and /2,
measured from the perigee direction and relating to the points of, respectively,
shadow entrance and shadow exit, by using another definition of the shadow
function than that defined above. Escobal’s shadow function, denoted by g(/), is
defined as follows
gð/Þ R2E ð1 þ e cos /Þ2 þ ½að1 e2 Þ2 ðsx cos / sy sin /Þ2 ½að1 e2 Þ2
where, again, sx and sy are the projections of the Sun–Earth unit vector (s) onto,
respectively, the x-axis and the y-axis of the geocentric perifocal system xyz. The
geometrical meaning of the expression
sx cos / sy sin /
is the scalar product (i.e. the cosine of the angle w) between the anti-Sun–Earth unit
vector (−s) and the position unit vector (r/r) connecting the centre of the Earth with
the satellite at the points of shadow entrance and shadow exit. The anti-Sun–Earth
unit vector (−s) is also the geocentric unit position vector of the Sun, that is,
rS
s ¼
jrS j
As Escobal [9, 118] points out, w must fall in the interval [90°, 270°], because
this is the only domain where a shadow exists. This is the same condition as that
indicated by Cunningham (that is, 90° w1 180° and 180° w2 270°).
In other words, the following relation
sx cos / sy sin /\0
must hold. Therefore, the criterion which identifies the two meaningful zeros (/1
and /2) of the shadow function g(/) is
sx cos / þ sy sin / [ 0
because, in any other case, the satellite would be in direct sunlight. Escobal also
points out that a change of the sign in front of g(/) from minus to plus indicates the
entrance of the satellite into the shadow of the Earth, whereas a change of sign from
plus to minus indicates the exit of the satellite from the shadow of the Earth.
An example of application of the methods described above is given below. The
subsequent section will show how to refine the results found by using the simple
mathematical models described above, in order to take account of the effects of the
perturbations. Let us consider a satellite revolving around the Earth according to the
following orbital elements
a ¼ 20000 km X ¼ 44 :602
e ¼ 0:43336 x ¼ 30 :683
i ¼ 30 :193 /0 ¼ 350 :85
at epoch t0 = 8 March 2012, noon UT1. As has been shown in section 14 of the
present chapter, the geocentric Cartesian co-ordinates of the Sun are
XS ¼ 0:972225757 AU
YS ¼ 0:184302887 AU
ZS ¼ 0:0798989791 AU
The corresponding co-ordinates of the Sun–Earth unit vector s are then (using
five significant figures)
XS
sX ¼ 1 ¼ 0:97932
XS þ YS2 þ ZS2 2
2
YS
sY ¼ 12 ¼ 0:18565
XS2 þ YS2 þ ZS2
ZS
sZ ¼ 12 ¼ 0:080482
XS2 þ YS2 þ ZS2
The Sun–Earth unit vector, in the geocentric perifocal reference system xyz, has
the following components
sx ¼ ðcos X cos x sin X sin x cos iÞsX þ ðsin X cos x þ cos X sin x cos iÞsY
þ ðsin x sin iÞsZ ¼ ðcos 44 :602 cos 30 :683 sin 44 :602 sin 30 :683
cos 30 :193Þ ð0:97932Þ þ ðsin 44 :602 cos 30 :683 þ cos 44 :602
sin 30 :683 cos 30 :193Þ 0:18565 þ ðsin 30 :683 sin 30 :193Þ 0:080482
¼ 0:10530
sy ¼ ð cos X sin x sin X cos x cos iÞsX þ ð sin X sin x þ cos X cos x cos iÞsY
þ ðcos x sin iÞsZ ¼ ð cos 44 :602 sin 30 :683 sin 44 :602 cos 30 :683
cos 30 :193Þ ð0:97932Þ þ ð sin 44 :602 sin 30 :683 þ cos 44 :602
cos 30 :683 cos 30 :193Þ 0:18565 þ ðcos 30 :683 sin 30 :193Þ 0:080482
¼ 0:93351
sz ¼ ðsin X sin iÞsX þ ð cos X sin iÞsY þ ðcos iÞsZ ¼ ðsin 44 :602 sin 30 :193Þ
ð0:97932Þ þ ð cos 44 :602 sin 30 :193Þ 0:18565 þ ðcos 30 :193Þ
0:080482 ¼ 0:34275
Hence, c and d are given by

2 !12 3 "
12 #
s2z 0:342752
d ¼ arctan4 5 ¼ arctan ¼ 0:34984 radians
s2x þ s2y 0:105302 þ 0:933512

¼ 20 :044

sy 0:93351
c ¼ arctan ¼ arctan ¼ 1:4585 radians ¼ 83 :564
sx 0:10530
Now we compute the upper limit (d′) of d as follows

0 RE ð1 e cos cÞ 6371 ð1 0:43336 cos 1:4585Þ
d ¼ arcsin ¼ arcsin
að1 e2 Þ 20000 ð1 0:433362 Þ
¼ 0:38241 radians ¼ 21 :910
Since d = 20°.044 is less than d′ = 21°.910, then the satellite spends a fraction
of its orbital period T in the shadow of the Earth.
We search now the two meaningful zeros w1 and w2, relating to the true anomalies
of, respectively, shadow entrance and shadow exit, of the shadow function
a 1
f ðwÞ 1 þ e cosðw cÞ 1 e2 1 cos2 w cos2 d 2
RE
where w1 falls inside the interval [p/2, p] and w2 falls inside the interval [p, 3p/2].
In order to determine w1, we evaluate f(w) at, respectively, the lower end (p/2)
and the upper end (p) of the interval, as follows
f ð1:5708Þ ¼ 1 þ 0:43336 cosð1:5708 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 1:5708 cos2 0:34984Þ1=2 ¼ 1:1190
f ð3:1416Þ ¼ 1 þ 0:43336 cosð3:1416 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 3:1416 cos2 0:34984Þ1=2 ¼ 0:077543
Since f(w) changes sign in the interval 1.5708 w 3.1416, then there is
assurance of a zero (w1) of the shadow function f(w) falling inside this interval. By
applying Müller’s method, we choose arbitrarily another value falling between
1.5708 and 3.1416. Let 2.8 be this value. At w = 2.8 radians, there results
f ð2:8000Þ ¼ 1 þ 0:43336 cosð2:8000 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 2:8000 cos2 0:34984Þ1=2 ¼ 0:087881
Now we set
x2 ¼ 1:5708 f2 f ðx2 Þ ¼ 1:1190

x0 ¼ 2:8000 f0 f ðx0 Þ ¼ 0:087881
x1 ¼ 3:1416 f1 f ðx1 Þ ¼ 0:077543
(where the subscripts used above have nothing to do with the points of shadow
entrance or exit or with epoch) and compute
h1 x1 x0 ¼ 3:1416 2:8000 ¼ 0:3416

h2 x0 x2 ¼ 2:8000 1:5708 ¼ 1:2292
c h2 =h1 ¼ 1:2292=0:3416 ¼ 3:5984
A ½cf1 f0 ð1 þ cÞ þ f2 =½ch21 ð1 þ cÞ ¼ ½3:5984 0:077543 þ 0:087881 ð1

þ 3:5984Þ 1:1190=½3:5984 0:34162 ð1 þ 3:5984Þ ¼ 0:22573
B ðf1 f0 Ah21 Þ=h1 ¼ ð0:077543 þ 0:087881 þ 0:22573 0:34162 Þ=0:3416
¼ 0:56137
C f0 ¼ 0:087881
where c has nothing to do with the angle formed by the line of apsides of the
satellite orbit and the projection of the Sun–Earth unit vector onto the orbital plane.
This done, a better approximation than x0 = 2.8 to the unknown value of w1 is
2C
x ¼ x0 1
B ðB2 4AC Þ2
where the upper sign (plus) holds if the value of B is greater than zero, whereas the
lower sign (minus) holds if the value of B is lower than zero. If B = 0, either sign
can be used. In the present case (B > 0), there results
2 0:087881
x ¼ 2:8 þ 1 ¼ 2:9679
0:56137 þ ð0:561372 4 0:22573 0:087881Þ2
The corresponding value of the function f(w) is
f ð2:9679Þ ¼ 1 þ 0:43336 cosð2:9679 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 2:9679 cos2 0:34984Þ1=2 ¼ 0:059613
Since 2.9679 falls between 2.8 and 3.1416, then we discard 1.5708 and choose
2.8, 2.9679 and 3.1416 for the next iteration. At the same time we reset the sub-
scripts, so as to have
x2 ¼ 2:8000 f2 f ðx2 Þ ¼ 0:087881

x0 ¼ 2:9679 f0 f ðx0 Þ ¼ 0:059613
x1 ¼ 3:1416 f1 f ðx1 Þ ¼ 0:077543
By so doing, we find, with five significant figures, w1 = 2.8875

radians = 165°.44.
Likewise, in order to determine w2, we evaluate f(w) at, respectively, the lower
and upper ends of the interval [p, 3p/2], as follows
f ð3:1416Þ ¼ 1 þ 0:43336 cosð3:1416 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 3:1416 cos2 0:34984Þ1=2 ¼ 0:077543
f ð4:7124Þ ¼ 1 þ 0:43336 cosð4:7124 1:4585Þ ð20000=6371Þ ð1 0:433362 Þ

ð1 cos2 4:7124 cos2 0:34984Þ1=2 ¼ 1:9803
Since f(w) changes sign in the interval 3.1416 w 4.7124, then there is
assurance of a zero (w2) of f(w) falling inside this interval.
By applying again Müller’s method, we find, with five significant figures,
w2 ¼ 3:2441 radians ¼ 185 :87
The corresponding true anomalies of the points of, respectively, shadow entrance
and shadow exit, measured from the perigee direction, are
/1 ¼ w1 c ¼ 2:8875 1:4585 ¼ 1:4290 radians ¼ 81 :876

/2 ¼ w2 c ¼ 3:2441 1:4585 ¼ 1:7856 radians ¼ 102 :31
Now, using Cunningham’s formula to compute the shadow time, t2 − t1, we

have
sin /1 sin 1:4290

p1 ¼ ¼ ¼ 0:93283
1 þ e cos /1 1 þ 0:43336 cos 1:4290
sin /2 sin 1:7856

p2 ¼ ¼ ¼ 1:0765
1 þ e cos /2 1 þ 0:43336 cos 1:7856
"
1
# "
1
#
1e 2 /1 1 0:43336 2 1:4290
q1 ¼ arctan tan ¼ arctan tan
1þe 2 1 þ 0:43336 2
¼ 0:49928
"
12
# "
1
#
1e /2 1 0:43336 2 1:7856
q2 ¼ arctan tan ¼ arctan tan
1þe 2 1 þ 0:43336 2
¼ 0:66287

1 i
200003 12
a3 2 h 1
t2 t1 ¼ e 1 e ð p1 p2 Þ þ 2ð q2 q1 Þ ¼
2 2
½0:43336
lE 398600:1
1
1 0:433362 2 ð0:93283 1:0765Þ þ 2 ð0:66287 0:49928Þ
¼ 1214:5 s ¼ 20m 15s
The same value can also be found by using Escobal’s formula, as shown below.
1 1
ð1 e2 Þ2 sin /1 ð1 0:433362 Þ2 sin 1:4290
sin Æ 1 ¼ ¼ ¼ 0:84069
1 þ e cos /1 1 þ 0:43336 cos 1:4290
e þ cos /1 0:43336 þ cos 1:4290

cos Æ 1 ¼ ¼ ¼ 0:54152
1 þ e cos /1 1 þ 0:43336 cos 1:4290
1 1
ð1 e2 Þ2 sin /2 ð1 0:433362 Þ2 sin 1:7856
sin Æ 2 ¼ ¼ ¼ 0:97012
1 þ e cos /2 1 þ 0:43336 cos 1:7856
e þ cos /2 0:43336 þ cos 1:7856

cos Æ 2 ¼ ¼ ¼ 0:24262
1 þ e cos /2 1 þ 0:43336 cos 1:7856

12

a3 sin Æ 2 sin Æ 1
t2 t1 ¼ arctan arctan eðsin Æ 2 sin Æ 1 Þ
lE cos Æ 2 cos Æ 1
By substituting the values computed above into the preceding equation, we

obtain
t2 t1 ¼ 1214:5 seconds
Finally, using Lambert’s formula, we compute the radii vectores r1 and r2,
joining the centre of the Earth with the points of, respectively, shadow entrance and
shadow exit, as follows
að 1 e 2 Þ 20000 ð1 0:433362 Þ
r1 ¼ ¼ ¼ 15307 km
1 þ e cos /1 1 þ 0:43336 cos 1:4290
að 1 e 2 Þ 20000 ð1 0:433362 Þ
r2 ¼ ¼ ¼ 17897 km
1 þ e cos /2 1 þ 0:43336 cos 1:7856
The chord joining the two points indicated above results from
1
c ¼ r12 þ r22 2r1 r2 cosð/2 /1 Þ 2 ¼ ½153072 þ 178972 2 15307 17897
1
cosð1:7656 1:4290Þ2 ¼ 6417:1 km
Lambert’s angles a and b are computed as follows
a r þ r þ c12
15307 þ 17897 þ 6417:112
1 2
sin ¼ ¼ ¼ 0:70375
2 4a 4 20000

1
r 1 þ r 2 c 2 15307 þ 17897 6417:1 2
1
b
sin ¼ ¼ ¼ 0:57865
2 4a 4 20000
Hence
h ai
a ¼ 2arcsin sin ¼ 2 arcsinð0:70375Þ ¼ 1:5613 radians
2

b
b ¼ 2arcsin sin ¼ 2 arcsinð0:57865Þ ¼ 1:2341 radians
2
Finally, the time t2 − t1 spent by the satellite in the shadow of the Earth is

12
12
a3 200003
t2 t1 ¼ ½ða sin aÞ ðb sin bÞ ¼ ½ð1:5613
lE 398600
sin 1:5613Þ ð1:2341 sin 1:2341Þ ¼ 1214:5 s
Escobal’s method searches the zeros (/1 and /2, relating to the true anomalies
of, respectively, shadow entrance and shadow exit) of the following shadow
function
Remembering that the condition of shadow entrance and exit is g(/) = 0 and
that the radius vector is r = a(1 − e2)/(1 + e cos /), the condition g(/) = 0 can also
be written as follows

2
RE 2
þ sx cos / sy sin / 1 ¼ 0
r
Using Escobal’s criterion to search only the two meaningful zeros, we impose
that
sx cos / sy sin /\0
Using the results found above by means of Cunningham’s method, we search /1

in the vicinity of 81°.876 and /2 in the vicinity of 102°.31. By using the values
computed above for of sx, sy, r1, cos /1, sin /1, r2, cos /2, and sin /2, it is easy to
verify that the condition g(/) = 0 holds with /1 = 1.4290 and /2 = 1.7856 radians.
The following part of the present section will show how to correct the values of
the approximate true anomalies /1 and /2, computed above, in order to take the
perturbations into account.
As to the Earth flattening, the Earth is represented now as an ellipsoid of rota-
tion, that is, as an oblate spheroid having a minor semi-axis Rpol (at the poles) and a
major semi-axis Req (at the equator), as shown in the following figure.
Let
Req Rpol
f ¼
Req
be the flattening of the Earth. Following Escobal [9, 118], the component along the
Earth polar axis of the Earth radius R, at the moment of shadow entrance or exit, is
að 1 e 2 Þ
RZ ¼ sin x sini cos / þ cos x sini sin / sz sx cos / þ sy sin /
1 þ e cos /
where /, the true anomaly of the satellite measured from the perigee direction, has
nothing to do with the geodetic latitude denoted by u in the preceding figure. Since
RZ can also be computed at any time as follows
1
Req ½1 ð2f f 2 Þ2 sin u
RZ ¼ 1
½1 ð2f f 2 Þ cos2 u 2
where u* (shown in the preceding figure) is the geocentric latitude of the point P
where the shadow ray departs, then the preceding expression, solved for cos2u*,
yields
R2eq ½1 ð2f f 2 Þ R2Z

cos2 u ¼
R2eq ½1 ð2f f 2 Þ R2Z ð2f f 2 Þ
Since an approximate value of RZ is known, the corresponding approximate

value of cos2u* can be computed by means of the preceding expression. This
approximate value of cos2u* makes it possible to compute a better approximation
to the value of the Earth radius, as follows
1
Req ½1 ð2f f 2 Þ2
R¼ 1
½1 ð2f f 2 Þ cos2 u 2
Now, this value of R is used instead of RE in the shadow function g(/)
and an iterative procedure is applied to find the new zeros (/1 and /2) of the
preceding equation. These new zeros are introduced into
að 1 e 2 Þ
RZj ¼ sin x sin i cos / þ cos x sin i sin / sz sx cos / þ sy sin /
1 þ e cos /
(where the subscript j indicates an iteration) and a new value of RZ is computed; this
value is introduced into
R2eq ½1 ð2f f 2 Þ R2Z

cos2 u j ¼
R2eq ½1 ð2f f 2 Þ R2Z ð2f f 2 Þ
and a new value of cos2u* is computed; this value makes it possible to get a better
approximation to the correct value of the Earth radius, as follows
1
Req ½1 ð2f f 2 Þ2
Rj ¼ 1
½1 ð2f f 2 Þ cos2 u 2
and so on, till the Earth radius does not vary any more within a chosen tolerance.
This iterative procedure makes it possible to take account of the effects of the Earth
flattening.
As to the correction for umbra and penumbra, due to the conical shape of the
umbra region, let us consider the semi-aperture angles fU and fP, shown below,
where the subscripts U and P indicate, respectively, the umbra and penumbra cones.
Following again Escobal [9, 118], the semi-aperture angles fU and fP are given
by
h i12
rS2 ðRS RE Þ2 RS RE
cos fU ¼ sin fU ¼
rS rS
h i12
rS2 ðRS þ RE Þ2 RS þ RE
cos fP ¼ sin fP ¼
rS rS
as shown in the following figure, where the Sun and the Earth are assumed to be
perfect spheres. However, this assumption does not affect the results which will be
shown below, because the Earth and Sun flattening has negligible effects on fU and
fP. Consequently, in order to take the umbra/penumbra effects into account,
Escobal [9, 118] redefines the shadow function g(/) as shown below:
gi ð/Þ R2Ei ð1 þ e cos /U Þ2 þ ½að1 e2 Þ2 ðsx cos /i sy sin /i Þ2 ½að1 e2 Þ2 cos2 fi
2að1 e2 ÞREi ðsx cos /i sy sin /i Þð1 þ e cos /i Þ sin fi
where the subscript i denotes umbra and penumbra distances, the upper sign is to be
taken at penumbra entrance and 0° fi 90°. Then, an iterative root-finding
procedure is applied to the shadow function gi(/) shown above.
As pointed out by Escobal, if the Earth flattening effects are also considered,
then the variables having an i subscript also have a j subscript, as shown in the case
of Earth flattening. In this case, the correction for Earth flattening is to be included
into the correction for umbra/penumbra.
As to the correction due to the Earth motion around the Sun, the quantities sx, sy,
RZ, and the angles fi which appear in the expressions given above vary, all of them,
with time, because their values depend upon the position of the Sun with respect to
the Earth.
Following again Escobal [9, 118], when the first-approximation values of the
shadow-entrance (/1) and shadow-exit (/2) true anomalies are known, then the
time (t) corresponding at a particular true anomaly (/) is

12
a3
t¼ ðÆ e sin Æ Þ þ s
lE
where s is the time of perigee passage, and Æ is the eccentric anomaly.

By so doing, the times t1 and t2 (corresponding to the first-approximation values
of, respectively, /1 and /2) can be computed. Now, a new computation of the
geocentric co-ordinates of the Sun (or an interpolation of the tables giving such
co-ordinates) yields the corresponding values of the quantities sx, sy, RZ, and fi.
As to the correction due to the atmospheric refraction, the rays of light coming
from the Sun have hitherto been assumed to travel along straight lines. In practice,
such rays are refracted during their passage through the Earth atmosphere, as shown
in the following figure, because they go through media having different densities.
The path resulting from this phenomenon satisfies Fermat’s principle, which
states that light follows the path of least time. Consequently, the umbra region is not
a circular cone, that is, a surface formed by a rotating straight line, but a surface
formed by a bent line following the path of a refracted ray. In other words, a ray of
light coming from the Sun at an angle fU with respect to the horizontal when
entering the atmosphere is deviated and leaves the atmosphere at an angle fU + DfU
with respect to the horizontal.
The angles fU (for umbra) and fP (for penumbra) can be expressed as a function
of the Sun radius, the Earth radius, and the Sun–Earth distance, as will be shown
below. They are incidence angles. Consequently, the problem reduces to evaluating
their variations DfU and DfP due to refraction. These values, once determined, can
be introduced into the expression shown above
gi ð/Þ R2Ei ð1 þ e cos /U Þ2 þ ½að1 e2 Þ2 ðsx cos /i sy sin /i Þ2 ½að1 e2 Þ2 cos2 fi
2að1 e2 ÞREi ðsx cos /i sy sin /i Þð1 þ e cos /i Þ sin fi
to take account of refraction. To this end, Dreher [113] makes the following
assumptions:
• the Earth atmosphere is formed by an infinite number of successive strata, each
of which has a uniform density;
• the successive strata are all concentric with the centre O of a spherical Earth of
radius RE;
• the value q0 of the atmospheric density q at the surface of the Earth (indicated
by the subscript 0) is q0 = 0.0027 slug/ft3 = 1.392 kg/m3;
• the value of the refractive index n at the surface of the Earth is n0 = 1.00029;
• the atmospheric density q decreases exponentially as the radius r increases, that
is, q(r) = e−r;
• the relation between the refractive index n and the atmospheric density q is
given by Sellmeier’s equation, that is, by n2(r) = 1 + [(n20 − 1)/q0]q(r);
• the upper limit of the atmosphere is h = 6419.719 km.
In practice, the atmosphere is not spherically stratified, because the Earth is an
oblate ellipsoid; nor does the density of each stratum remain constant, because of
movements of air and gradients of temperature and pressure. However, these per-
turbations can be neglected for many purposes.
To determine the refracted ray equation, Dreher uses Fermat’s principle,
according to which the path of a light ray passing through the points A and B of a
medium having a variable refractive index n is such that the following integral
ZB
‘¼ n ds
A
is an extremum, where ‘ is the optical path length, n is the refractive index, and ds is
the differential of arc length along the path AB.
The deviation angle has been determined by Dreher from a definite integral
which depends on the refraction index n = 1 + N 10−6, where N is refractivity.
The latter, in turn, is a function of wavelength and of the state of the atmosphere
(temperature, pressure, and humidity). The dependence of refractivity on humidity
is weak, as has been shown by Wittmann [120], and is usually neglected.
The refraction integral has been evaluated by Dreher, who has found the fol-
lowing value
DfU ¼ 380 :8
of the deviation angle, where the symbol ′ indicates minutes of arc.

Dreher does not specify the atmospheric conditions and the wavelength on which
his calculations are based. However, the value DfU = 38′.8 may be checked against
that resulting from Seidelmann’s formula [43] for altitudes between 0° and 15°:
34:133 þ 4:197h þ 0:00428h2

DfU ¼
1 þ 0:505h þ 0:0845h2
where h is the apparent angular altitude in degrees (in this case, h = 0°) and DfU is
expressed in minutes of arc. This formula, applied to the case of mean atmospheric
conditions (with a pressure of 101000 Pa, a temperature of 10 °C, and a wavelength
of 550 nm), yields the value DfU = 34′.133 for the deviation angle. Taking into
account variations of pressure p (Pa) and temperature t (°C), the previous value
(34′.133) must be multiplied by the approximate factor indicated by Meeus [40]
p 283
101000 273 þ t
Therefore, the angle of deviation at p = 760 mm of Hg (101325 Pa) and

t = 0 °C is

101325 283
DfU ¼ 340 :133 ¼ 350 :497
101000 273
which does not differ very much from the value (38′.8) computed by Dreher.
Garfinkel [121] and, more recently, Auer and Standish [122] have shown that the
problem of computing astronomical refraction for any value of the apparent angular
altitude can be reduced to a numerical quadrature, which can be performed by
using, for example, Simpson’s rule or the Gauss–Chebyshev method. This
quadrature becomes simple and exempt from singularities by choosing properly the
independent variable of integration. To this end, Auer and Standish choose the
angle between the radius vector and the light ray as the variable of integration, and
apply the quadrature method to a piecewise polytropic atmosphere, comprising two
layers, the first of which (troposphere), going from the surface of a spherical Earth
to a height ht (tropopause) of about 11 kilometres, is characterised by a constant and
finite temperature gradient dt/dr, whereas the second (stratosphere), going from the
tropopause to infinity, is isothermal. Air temperature and density are assumed to
vary continuously across the tropopause. Unlike Dreher [112], both Garfinkel [121]
and Auer and Standish [122] assume the validity of the Gladstone–Dale relation
n 1 ¼ cq
where n is the refraction index, q is the density, and c is a constant which depends
only on the substance (air) and the wavelength of light. The constant c has a finite
value within the troposphere and an infinite value in the stratosphere.
In case of a ray of light which barely touches the surface of the Earth (at an
apparent angular altitude of zero degrees, 0 °C and 760 mm of Hg), Garfinkel has
found a deviation angle of 2206″.4 = 36′.8. The deviation angle found by Auer and
Standish in the same conditions is 2189″.42 = 36′.5. Both of these values are nearer
to the value (35′.5) resulting from Seidelmann’s formula than is the value (38′.8)
computed by Dreher.
The other quantities considered above are fU and fP, that is, the semi-aperture
angles of, respectively, the umbra cone and the penumbra cone in the absence of
refraction. The value for fU has been computed geometrically by Dreher [113] and
found to be 15′. This value is substantially the same as that resulting from the
formulae given by Escobal [9, 118], that is,
h i12
rS2 ðRS RE Þ2 RS RE
cos fU ¼ sin fU ¼
rS rS
h i12
rS2 ðRS þ RE Þ2 RS þ RE
cos fP ¼ sin fP ¼
rS rS
which lead to fU = 15′.855 and fP = 16′.148.

Therefore, the post-refraction semi-aperture angle of the umbra cone results
fU þ DfU ¼ 150 :855 þ 350 :497 ¼ 510 :352
Likewise, the post-refraction semi-aperture angle of the penumbra cone results
DfP fP ¼ 350 :497 160 :148 ¼ 190 :349
3.18 The Radiation Pressure Due to the Earth
By Earth radiation pressure we mean the visible and thermal radiation emitted and
reflected by the Earth. This radiation is due to the reflected visible (wavelength in
the range 0.4–0.7 lm) and reradiated infrared (in the range 3.55–3.90 and 10.50–
12.50 lm, according to ref. [123]) energy hitting the spacecraft surface.
The only energy source of practical importance for the Earth is the Sun. Over the
whole surface of the Earth, about 30% [124] of incoming solar energy is reflected
back to space. According to Adhya [125], when the incoming energy into the
illuminated hemisphere is averaged over the whole surface of the Earth, the mean
incident energy flux is about 350 W/m2, of which approximately one third is
reflected in the visible part of the spectrum (by clouds, continents, and oceans), and
the rest is absorbed (by the atmosphere, soil, and water). The absorbed energy is
re-emitted by the Earth as infrared radiation. The short-wave Earth radiation is
directly reflected, in both the diffuse and specular mode, and is therefore present
only in the illuminated hemisphere of the Earth. By contrast, the long-wave Earth
radiation is always present to some extent even in the dark hemisphere of the Earth
[126].
The earliest mathematical models determined the Earth radiation pressure acting
on a satellite by computing a portion of solar radiation pressure for a given
Sun-Earth-satellite (w) angle. The direction was assumed to be purely radial.
Following Rodriguez-Solano [127], let us consider the geocentric reference
system xyz, illustrated in the following figure, and defined below.
3.18 The Radiation Pressure Due to the Earth 479
The fundamental (xy) plane contains the satellite, Earth, and Sun, all of which
are considered as point-masses. The x-axis coincides with the satellite position
vector r. The y-axis is orthogonal to r and lies in the satellite-Earth-Sun plane. The
z-axis is orthogonal to this plane, so as to form a right-handed system with x and
y. The unit vectors of the x, y, and z axes defined above are indicated with ux, uy,
and uz, respectively. Let w be the satellite-Earth-Sun angle. Let v be angle between
the diffusely reflected radiance of the Earth and the z-axis. Let u be the angle
between the satellite position vector r and the projection of the reflected radius onto
the xy-plane. We consider now another set of unit vectors ur, us, and un, whose
components in the xyz system are shown below
2 3 2 3 2 3
1 cos w sin v cos u
ur 4 0 5 us 4 sin w 5 un 4 sin v sin u 5
0 0 cos w
In other words, ur coincides with ux, us points towards the Sun, and un points in
the direction of the diffusely reflected radius. The satellite-Earth-Sun angle (w) is
also illustrated in the following figure, where the angle b0 is the elevation of the Sun
above the orbital plane of the satellite.
The mathematical model described by Rodriguez-Solano [127] is based on the

following assumptions:
• the Earth is a sphere which behaves as a Lambertian reflector (by the way, a
Lambertian surface is an ideal diffusely reflecting surface, which reflects the
incident solar radiation isotropically, according to Lambert’s cosine law);
• the value of the total terrestrial albedo (the ratio of reflected radiant power to
incident radiant power, at all wavelengths and in all directions) is 0.3, in
accordance with Goode et al. [124];
• there is conservation of energy as a whole, that is, all the energy coming from
the Sun and hitting the Earth must go back to space as either reflected or emitted
radiation; and
• the average radius of the Earth is used in the model.
In order to obtain an analytical solution, two further assumptions are made:
• the distance, d, between the radiating surface elements and the satellite is
constant; and
• the vector, r, of the satellite from a given surface element is also constant.
Let c and h be, respectively, the angle of incidence of the incoming radiation and
the angle of reflection of the outgoing radiation with respect to the normal, un, to a
surface element. These angles are shown in the following figure.
The incident radiant power, Pin, on a planar surface, A, has the following
expression
Pin ¼ Ssc A cos c
where Ssc 1367 W/m2 is the solar constant. Let a 0.3 and d be, respectively,
the terrestrial albedo and the distance of the satellite from the surface A.
Since A is assumed to be a Lambertian surface, then the reflected irradiance,
which comes from the Earth and hits the satellite at a distance d from A, is
a a
Erefl ¼ Pin cos h ¼ ðSsc A cos cÞ 2 cos h
pd 2 pd
The preceding expression holds with cos h

0 and cos c
0.
Now, let us consider the Earth as a Lambertian sphere, made up of an infinite
number of planar elementary surfaces dA. The cosines of the angles of incidence (c)
and reflection (h) can be expressed as follows
cos h ¼ ur un ¼ sin v cos u

cos c ¼ us un ¼ sin v cosðw uÞ
where u, v, and w are the angles defined above. Introducing these expressions into
a
dE refl ¼ ðSsc dA cos cÞ cos h
pd 2
yields
a
dE refl ¼ 2
Ssc sin2 v cosðw uÞcos u dA
pd
which is the reflected irradiance coming from a single elementary surface dA of the
Earth to the satellite. The preceding expression, integrated over the illuminated part
of the spherical Earth visible to the satellite, with
dA ¼ R2E sin v dv du
0vp
w p=2 u p=2
0wp
RE being the mean radius of the Earth, yields the reflected irradiance hitting the
satellite as a function of w:
Zp Zp=2
a
Esat ðwÞ ¼ 2
Ssc sin2 v cosðw uÞcos u R2E sin v dv du
pd
0 wp=2
The integral indicated above can be evaluated separately for the two variables v
and u, as follows
Z Z Z

sin3 v dv ¼ sin2 v ðsin vÞ dv ¼ sin2 vð cos vÞ ð cos vÞ2 sin v cos v dv
Z Z
¼ sin2 v cos v þ 2 cos2 v sin v dv ¼ sin2 v cos v 2 cos2 v dðcos vÞ
¼ sin2 v cos v 2ð1=3 cos3 vÞ ¼ sin2 v cos v 2=3 cos3 v

¼ ð1 cos2 vÞ cos v 2=3 cos3 v ¼ cos v þ cos3 v 2=3 cos3 v
¼ 1=3 cos3 v cos v
This integral, over the interval [0, p], yields

1 3 1 3 1 1 4
cos p cos p cos 0 cos 0 ¼ þ 1 þ 1 ¼
3 3 3 3 3
Z Z
cosðw uÞ cos u du ¼ ðcos w cos u þ sin w sin uÞ cos u du
Z Z
¼ cos w cos2 u du þ sin w sin u cos u du
Z Z Z
cos u du ¼ cos u sin u
2
sin uð sin u duÞ ¼ cos u sin u þ
sin2 u du
Z Z
¼ cos u sin u þ ð1 cos2 uÞdu ¼ cos u sin u þ u cos2 u du
¼ 1=2ðu þ cos u sin uÞ

Z Z
sin u cos u du ¼ sin u dðsin uÞ ¼ 1=2 sin2 u
R
The integral cosðw uÞ cos u du, over the interval [w − p/2, p/2], yields
1 1 1
cos wðp w þ sin w cos wÞ þ sin w 1 cos2 w ¼ ½ðp wÞcos w þ sin w
2 2 2
Therefore, the reflected irradiance hitting the satellite is
a 41
Esat ðwÞ ¼ S R2
2 sc E 3 2
½ðp wÞcos w þ sin w
pd
2 a
¼ Ssc R2E ½ðp wÞcos w þ sin w
3 pd 2
By setting AE = p R2E, the preceding expression becomes
2 a
Esat ðwÞ ¼ Ssc AE ½ðp wÞcos w þ sin w
3 ðpd Þ2
This is the total irradiance due to the albedo of the Earth as a function of the
satellite-Earth-Sun angle, w, and the Earth–satellite distance, d, which has so far
been considered constant. This distance includes the mean radius of the Earth, RE,
and the satellite altitude, h, from the Earth surface, so that
d ¼ RE þ h
Rodriguez-Solano [127] points out that the preceding expression holds only
when the altitude of the satellite is much greater than the mean radius of the Earth
(i.e. h RE). When this condition is satisfied, the irradiance vector Esat(w, h) due
to the total albedo of the Earth is directed radially, and the magnitude Esat(w, h) of
this vector depends only on w and h, as follows
( )
2 a
Esat ðw; hÞ ¼ Ssc AE ½ðp wÞcos w þ sin w ur
3 ½pðRE þ hÞ2
By contrast, when the actual direction and distance from a given surface element
to the satellite are used, then the irradiance vector is computed numerically, as will
be shown below. In this case, with reference to the following figure, the reflected
irradiance coming from a single elementary surface dA of the Earth to the satellite is

a
dErefl ¼ ðSsc dA cos cÞ cos h ud
pd 2
and is directed along the unit vector ud. Such is the case only if cos h
0 and
cos c
0. In any other case, dErefl is zero.
As has been shown above, the cosines of the angles c (relating to the incident
radiation coming from the Sun) and h (relating to the radiation reflected by the
Earth) are
cos h ¼ ud un
cos c ¼ us un
The vector d goes from the elementary surface element dA to the satellite; the
vector r goes from the geocentre to the satellite; and the vector s goes from the
geocentre to the Sun. The magnitude of r is |r| = RE + h. The unit vectors ud, ur,
and us are such that ud = d/|d|, ur = r/|r|, and us = s/|s|. The unit vector urp indicates
the direction perpendicular to r. When the irradiance vector is computed numeri-
cally, then the vectors d and r have two distinct directions, as shown in the pre-
ceding figure. By contrast, when the irradiance vector is computed analytically, then
d and r have one and the same direction. The angle b indicates the part of the Earth
which can be seen by the satellite. In the particular case of a GPS satellite (whose
altitude is h 20200 km), the angle b results

RE 6371
b ¼ arccos ¼ arccos ¼ 76
RE þ h 6371 þ 20200
When the value of b is known, it is possible to compute the area of the spherical
surface seen by the satellite, as follows
1
A ¼ 2pR2E ð1 cos bÞ ¼ AT ð1 cos bÞ
2
where AT = 4pR2E is the area of the whole Earth. Therefore, the area of the surface
seen by a GPS satellite is
1 1
A ¼ ð1 cos bÞAT ¼ ð1 cos 76 ÞAT ¼ 0:379AT
2 2
that is, about 38% of the area of the whole Earth.

Now, in order to integrate dErefl = (Ssc dA cos c) [a/(pd2)] cos h, we replace the
elementary area dA with R2E sin v dv du. Therefore, the integration limits are
0vp
w p=2 u p=2
0wp
The integral of dErefl over the illuminate part of the Earth visible to the satellite is
a vector having two components. They are:
R
(1) a radial component [Esat(w, h)]r = dErefl ur , and
R
(2) a perpendicular component [Esat(w, h)]rp = dErefl urp :
There is no component of Erefl out of the satellite-Earth-Sun plane, because the
illumination is symmetric with respect to the axis perpendicular to this plane for
constant albedo.
The two components (radial and perpendicular) of the irradiance vector, Erefl,
cannot any more be computed analytically. When these components have been
computed numerically, the total radiant power reflected results from integrating the
radial component, [Esat(w, h)]r, of the irradiance vector over the surface, A, of the
sphere of radius rsat = RE + h, as follows
Z
Prefltotal ¼ ½Esat ðw; hÞr dA ¼ aPinctotal
A
where dA = 2p(RE + h)2 sin w dw (with 0 w p) is the elementary surface, a

is the terrestrial albedo, and Pinc-total = p R2E Ssc is the total incident radiant power
coming from the Sun and intersected by the Earth.
As has been shown at the beginning of the present section, a fraction of this
power is reflected back by the Earth in the visible spectrum; the other fraction is
absorbed and then re-emitted as infrared radiation over the total surface (4pR2E) of
the Earth.
Since a = 0.3 is the reflected part of the incoming power, the power per square
metre radiated by the Earth at its surface can be expressed as follows
1
EEarth ¼ ð1 aÞSsc
4
where Ssc 1367 W/m2 is the solar constant. The following part of the present
section will consider not the visible radiation reflected back, which has been taken
in consideration above, but rather the infrared radiation re-emitted by the Earth.
Here, too, the irradiance of the Earth in the infrared spectrum can be computed
analytically or numerically. In the first case (analytical computation), it is necessary
to make the same assumptions as those indicated above.
The emitted irradiance, coming from an elementary surface dA of the Earth and
hitting the satellite at a distance d, is
1a
dEemit ¼ ðSsc cos hÞdA
4pd 2
where cos h = ur · un = sin v cos u and dA = R2E sin v dv du. As has been shown
above, RE = 6371 km is the mean radius of the Earth, v is the angle between the
diffusely reflected radiance of the Earth and the normal to the satellite-Earth-Sun
plane, and u is the angle between the geocentric position vector r of the satellite and
the projection of the reflected radius onto the satellite-Earth–Sun plane.
The emitted irradiance dEemit, integrated over the part of the spherical Earth
visible to the satellite at infinity, yields
Z Zp Zp=2
1a
Esat ðwÞ ¼ dEemit ¼ ½Ssc sin v cos uR2E sin v dv du
4pd 2
0 p=2
1a
¼ AE Ssc
4pd 2
where the notation Esat(w) is used to indicate that the emitted irradiance is in a
specific direction. However, Esat(w) is constant and does not depend on w.
The integral of Esat(w), over the spherical surface A of radius RE + h, is
Z Zp
1a
Pemittotal ¼ Esat ðwÞdA ¼ 2p AE Ssc ðRE þ hÞ2 sin w dw
4pd 2
A 0

RE þ h 2
¼ ð1 aÞAE Ssc ¼ ð1 aÞPinctotal
d
which shows that d must be equal to RE + h, as is also the case with the reflected
radiation. The irradiance vector at the satellite altitude, h, due to the radiation
emitted by the Earth, is then
" #
1 a AE Ssc
Esat ðw; hÞ ¼ ur
4p ðRE þ hÞ2
which holds only when h is much greater than RE.

In the second case (numerical computation), it is not necessary to assume a
constant distance, d, and a constant vector, r. The emitted irradiance vector, dEemit,
from a surface element towards the satellite is

1a
dEemit ¼ ðSsc cos hÞdA ud
4pd 2
This holds if cos h

0. Otherwise, dEemit is equal to zero.
The integral of dEemit over the part of sphere visible to the satellite is

Z
½Esat ðw; hÞr ¼ dEemit ur
for constant a, and the total power emitted is

Z
Pemittotal ¼ ½Esat ðw; hÞr dA ¼ ð1 aÞPinctotal
A
The mathematical model shown above has considered separately the two com-
ponents (reflected and emitted) of the radiation. The complete model results from
the addition of these components.
The Earth irradiance vector, E(w, h), acting on the satellite, if computed ana-
lytically (subscript A), results from

AE Ssc 2a 1a
EA ðw; hÞ ¼ ½ ðp w Þcos w þ sin w þ ur
ðRE þ hÞ2 3p2 4p
and has only the radial component.

The same vector, if computed numerically (subscript N), results from

Z
Z
½EN ðw; hÞr ¼ dErefl ur þ dEemit ur

Z
½EN ðw; hÞrp ¼ dErefl urp
and has two (radial and perpendicular) components.

In the preceding expressions, dErefl is equal to

a
ðSsc dA cos cÞ cos h ud
pd 2
only when cos h

0 and cos c
0. In any other case, dErefl is equal to zero.
Likewise, in the same expressions, dEemit is equal to

1a
ð S sc cos h ÞdA ud
4pd 2
only when cos h

0. Otherwise, dEemit is equal to zero.
The elementary radiating area, dA, is equal to
dA ¼ R2E sin v dv du
where
0vp
w p=2 u p=2
0wp
3.19 The Atmospheric Drag
The perturbation due to atmospheric drag affects all satellites, from those orbiting
the Earth at low altitudes (h 200 km) to those at and beyond the geostationary
altitude [128], the difference between the two cases residing in the magnitude of this
force. In practice, atmospheric drag is considered if the altitude of perigee is less
than 1000 km. In case of low-altitude satellites, atmospheric drag is the largest
perturbation among those of non-gravitational origin [2].
In general terms, drag is the resistance acting on a vehicle which moves through
a continuous medium. As is well known, all substances are essentially discontin-
uous, because they consist of molecules, which in turn are made of atoms.
However, some physical phenomena occur at length- and time-scales which are, by
several orders of magnitude, higher than those which are proper to the atomic
structure of matter.
The relative distances which separate atoms are of the order of a few ångströms
(1 å = 1 10−10 m), and the characteristic times of atomic bond vibrations are of
the order of femtoseconds (1 femtosecond = 1 10−15 s). By continuous medium
3.19 The Atmospheric Drag 489
we mean a finite amount of matter, whose physical properties do not depend on its
actual size or on the time interval over which they are measured.
In classical fluid dynamics, the mean free-molecule path (i.e. the mean value of
the distances travelled by each of the neutral particles, in a selected volume,
between successive collisions with other particles in that volume) is as short as
about 68 nm = 6.8 10−8 m at sea level and at a temperature of 25 °C [129]. By
contrast, there are cases in which the continuous medium scheme is not suitable to
describe the behaviour of the medium through which a vehicle moves. At altitudes
where artificial satellites revolve about the Earth, the atmospheric density is so low
that the flow around a satellite cannot be considered as continuous. In other words,
the atmosphere at such altitudes can no longer be considered a continuous medium,
but rather a plurality of discrete particles. This is because, in the upper atmosphere,
the mean free-molecule path, k, increases, with a roughly exponential law, to values
of over 200 metres at an altitude of 200 km [130]. Schlatter [131] indicates k = 240
m at this altitude (h = 200 km). When the mean free-molecule path becomes
comparable with a typical linear dimension, d, of the vehicle, then the properties of
the flow field depend upon the Knudsen number
k
K¼
d
When K is much greater than unity, the type (elastic or inelastic) of collision
between gas molecules and solid surfaces becomes of paramount importance, as
will be shown below.
The present section is concerned with the aerodynamic resistance to the motion
of a space vehicle (such as an artificial satellite, a sounding rocket, or an inter-
planetary spacecraft) travelling through the Earth atmosphere. In addition to the
aerodynamic drag, which acts constantly in a direction opposite to the motion of the
space vehicle, there is also an aerodynamic lift, which acts perpendicularly to the
direction of the vehicle motion, and whose value depends on the vehicle orienta-
tion. In general, the centre of pressure of the vehicle does not coincide with its
centre of mass. Therefore, the vehicle is also subject to a turning moment, M, about
its centre of mass. In case of controlled satellites, the turning moment due to the
atmosphere is counterbalanced by control jets, and consequently the satellite pro-
ceeds along its orbit without rotating.
In case of uncontrolled space objects, such as, for example, cylindrical-shaped
final-stage rockets, the turning moment causes such objects to tumble end-over-end,
with a continuous change of the lift direction. Therefore, the effects of lift tend to
cancel out and are negligible [132]. The drag and lift forces (respectively, fD and fL)
and the tumbling moment (M) are illustrated in the following figure.
The main equation which expresses aerodynamic drag is

1
jf D j ¼ CD p A ¼ CD qjvrel j2 A
2
where |fD| is the magnitude of the drag force vector fD, CD is a dimensionless
quantity called drag coefficient, p = ½ q |vrel|2 is the dynamic pressure on the
moving vehicle, q is the atmospheric density, |vrel| is the magnitude of the velocity
vector of the vehicle with respect to the atmosphere, and A is the projected area of
the vehicle normal to the relative velocity vector, vrel. This vector, in turn, is the
velocity of the incident air stream relative to the vehicle. Assuming that the
atmosphere, at least at the satellite altitudes, rotates with the Earth as a rigid body,
the relative velocity vector can be expressed as follows
vrel ¼ v vatm
where v is the velocity of the satellite with respect to the inertial reference system,
and vatm is the velocity vector of the Earth atmosphere with respect to the same
reference system. As shown above, the drag force vector, fD, is directed oppositely
to the relative velocity vector. This fact is expressed as follows
1
f D ¼ CD q A jvrel j vrel
2
Consequently, the drag acceleration vector acting on the satellite is

fD 1 CD A
aD ¼ ¼ q jvrel j vrel
m 2 m
where m is the mass of the vehicle. The quantity between parentheses, that is, the
drag coefficient times the area-to-mass ratio, is called the ballistic coefficient
CD A
B¼
m
and is measured in m2/kg. Some authors (Keesee [133], for one) define, as the
ballistic coefficient, another quantity, b, which is the inverse of the quantity
indicated above, that is, b = 1/B = m/(CDA). Gaposchkin and Coster [128] call the
drag coefficient, ballistic coefficient. In accordance with Sentman [134], the basic
assumptions of free-molecule flow are the following:
(1) the collisions between the incident molecules and the vehicle surface are much
more numerous than the collisions involving only the incident molecules; and
(2) the collisions between the incident molecules and the vehicle surface are much
more numerous than the collisions between the incident molecules and the
re-emitted molecules, that is, the incident flow is not disturbed by the presence
of the vehicle, and the equilibrium velocity distribution of the incident mole-
cules is changed only by collision with the vehicle.
According to Cook [130], in order for a free-molecule flow to occur for a body
moving at very high speed, the condition K 1 is not sufficient, if the surface
temperature is low and appreciable energy accommodation takes place.
The further condition to be satisfied for a free-molecule flow to exist is
vinc
K
vreem
where vinc is the velocity at which the molecules impinge on the body, and vre-em is
the velocity at which the molecules are re-emitted by the body. The largest value of
vinc/vre-em for a satellite is about 13, but this will only occur when the molecules are
re-emitted with a kinetic energy corresponding to the surface temperature [130].
The degree of energy transfer between a gas molecule and a surface is usually
expressed in terms of the accommodation coefficient, a, which is the ratio of the
energy change experienced by the impinging molecules to the maximum energy
change that could take place. Let Einc and Ere-em be the average kinetic energies of,
respectively, the incident molecules and the re-emitted molecules. Let Esurf be the
average kinetic energy that the re-emitted molecules would have if they left the
surface at the temperature of the surface. The accommodation coefficient a is
defined as follows
Einc Ereem
a¼
Einc Esurf
Roberts [135] showed that freshly cleaned surfaces of Tungsten and Nickel
stricken by atoms of Helium have low accommodation coefficients (respectively,
0.05–0.07 for Tungsten and 0.08 for Nickel), but also that these values become
considerably higher (respectively, 0.19 for Tungsten and 0.20 for Nickel) after
several hours, because contaminants collect on the surfaces, even in high-vacuum
experimental conditions. This is because the incident molecules are re-emitted by
contaminated surfaces with high accommodation coefficients in a diffuse angular
distribution. By contrast, the incident molecules are re-emitted by clean surfaces
with low accommodation coefficients in quasi-specular angular distribution.
The two types (specular and diffuse) of re-emission are illustrated in the fol-
lowing figure. Following Sentman [134], in the specular re-emission (shown on the
left-hand side of the figure), a gas molecule strikes the surface of the satellite at an
angle of incidence hinc and is then re-emitted, like a billiard ball, at an angle hre-em
equal to hinc. The re-emission velocity vre-em is not necessarily equal to the inci-
dence velocity vinc. By contrast, in the diffuse re-emission (shown on the right-hand
side of the figure), a gas molecule strikes the surface at an angle of incidence hinc
and is then re-emitted in a random direction, which means that there are as many
molecules with a particular velocity leaving the surface in one direction as there are
leaving the surface in the opposite direction. The angle of re-emission hre-em is by
no means related to the angle of incidence hinc. Again, the re-emission velocity
vre-em is not necessarily equal to the incidence velocity vinc.
In both types of re-emission, the velocity of a molecule re-emitted by the surface

depends on the amount of energy which the molecule transfers to the surface before
being re-emitted.
If Ere-em = Einc, then there is no exchange of energy between the incident
molecules and the surface. In this case, a = 0. By contrast, if Ere-em = Esurf, then the
incident molecules and the surface reach thermal equilibrium before the molecules
are re-emitted. In this case, a = 1.
In the definition of a, all energies associated with those molecular degrees of
freedom which enter into an energy exchange with the surface are implicitly
assumed to be accommodated to the same degree. Experimental evidence shows
that this holds for the translational and rotational degrees of freedom, whereas the
vibrational degrees of freedom are practically not affected by a surface collision
[134].
At the early times of the study of rarefied gas dynamics, the concepts shown
above have provided the basis for the purpose of determining the aerodynamic
forces and moments acting on a satellite orbiting the Earth in the upper atmosphere.
Successively, with the advent of the space era (the first artificial satellite, called
Sputnik 1, was injected into an elliptical Earth orbit on 4 October 1957), further
information on gas–surface interactions has come from measurements performed on
the orbital decay of artificial satellites. The surfaces of such satellites are covered
with adsorbed molecules (i.e. with atomic oxygen particles trapped on the surfaces
of the satellites), which affect the energy accommodation and angular distributions
of molecules leaving these surfaces in the range of altitudes going from 150 to
300 km [136].
The better understanding of these physical phenomena, in turn, has made it
possible to calculate the values of the drag coefficient, CD, as will be shown below.
The expression indicated above (½CDqA|vrel|2/m) for the magnitude of the drag
acceleration vector contains various quantities, whose values are known with dif-
ferent degrees of accuracy. The mass of the satellite is known within 1%. The
projected area, A, of the satellite normal to the incident velocity vector depends on
the satellite shape, size, and angle of attack to the incident flow. The attitude
stabilisation of the satellite controls the satellite orientation with respect to the
Earth-fixed velocity vector, and hence the projected area.
The atmospheric density, the relative velocity, and the drag coefficient are
known much less accurately than the satellite mass and projected area. The
atmospheric density and the relative velocity will be considered separately in the
next sections. The following part of the present section is concerned with the drag
coefficient.
As a result of what has been shown above, the value of the drag coefficient, CD,
depends on the behaviour of the gas particles after striking the surface of the
satellite, that is, on their velocity of re-emission. The gas flow around the satellite
can be considered according to the Newtonian impact theory: the gas particles, after
striking the satellite, lose their component of momentum normal to the satellite
surface, whereas their tangential component of momentum is conserved.
Following Keesee [133], let vinc be the velocity of an incident gas particle and let
vre-em be its velocity after re-emission. According to the second principle of dynamics,
the force of the incident particles is equal to the time derivative of their momentum
dðmvÞ
f ¼
dt
where m = q A vinc dt is the mass of an incident particle, A is the projected area of

the satellite normal to the incident velocity vector, vinc, and q is the air density.
According to the third principle of dynamics, the drag due to the incident particles
is equal and directed oppositely to f, that is,
dðmvÞ vreem vinc

fD ¼ f ¼ ¼ m
dt dt
By definition, the drag coefficient is
q A vinc dtðvreem vinc Þ

fD dt vreem vinc
CD ¼ 1 2 ¼ 1 2
¼ 2
2 q vinc A 2 q vinc A
vinc
If vre-em = 0 (which happens when all the incident molecules are diffusely
re-emitted), then CD = 2. By contrast, if vre-em = –vinc (which happens when all the
incident molecules are specularly re-emitted), then CD = 4. In practice, this fact
poses bounds to the range of variation of the drag coefficient.
The knowledge of this coefficient is intimately connected with that of air density.
According to Mance [137], three types of drag coefficients can be considered: fixed
drag coefficients, fitted drag coefficients, and physical drag coefficients.
A fixed drag coefficient of 2.2 was determined by Cook [130] in laboratory tests
carried out in 1965. In such tests, related to drag in hyperthermal free-molecule
flow, Cook found the results shown in the following table, for, respectively, a flat
plate normal to the flow, a flat plate at an incidence angle h with respect to the flow,
a cylinder perpendicular to the flow, and a cone of semi-vertex angle w with vertex
forwards and axis parallel to the flow. The quantity r appearing in the table denotes
the ratio vre-em/vinc.
The drag coefficient (based on the projected area perpendicular to the direction
of motion) was determined separately in the two cases of, respectively, diffuse
re-emission and accommodated specular re-emission. Cook found that, in case of a
sphere, the value CD = 2.2 corresponded to the value a 0.95 of the accommo-
dation coefficient.
Shape Drag coefficient

Diffuse re-em Accomm. specular re-em
Flat plate (normal to flow) 2(1 + 2
/3 r) 2(1 + r)
Flat plate at incidence h 2(1 + 2
/3 r sin h) 2[1 − r cos(2h)]
4
Sphere 2(1 + /9 r) 2
Cylinder perpendicular to flow 2(1 + 1
/6 pr) 2(1 + 1=3r)
Cone of semi-vertex angle w 2(1 + 2
/3 r sin w) 2[1 − r cos(2w)]
Successive measurements carried out in orbit by means of pressure (density)

gauges and mass spectrometers made it possible to calculate the drag coefficients of
satellites in low-Earth orbit up to 300 km with considerable confidence at times of
low-to-moderate solar activity [136]. In particular, with reference to four satellites
of compact shape, such measurements showed how much the drag coefficients of
such satellites can deviate from the value of 2.2, which has often been used in
deducing neutral densities from drag data [136].
Later on, at times of solar maximum, further measurements were carried out by
Pardini et al. [138], because at such times the greatest amount of atomic oxygen is
adsorbed on satellite surfaces, and consequently the greatest part of diffuse
re-emission was expected. Ten satellites were considered in the study cited above,
with perigee altitudes between 200 and 630 km. The orbital decay of such satellites
was analysed during the sunspot maxima of solar cycles 22 (1989–1990) and 23
(1999–2002). For each satellite, a series of drag coefficient values was obtained as a
function of the accommodation coefficient. The drag coefficients so determined are
called “physical drag coefficients”, in order to distinguish them from drag
coefficients determined by fitting the orbital decay to a particular model of atmo-

spheric density. The drag coefficients of the latter type are called “fitted drag
coefficients” and were computed by fitting, in the least-squares sense, the decay of
the major semi-axis. The fitted drag coefficients were then adjusted upward by a
certain amount, dependent on altitude, to account for the known biases in the
atmospheric density model used at sunspot maximum. The modified fitted drag
coefficients were called “observed drag coefficients” by Pardini et al. [138]. After
finding the physical drag coefficient matching the observed drag coefficient, the
corresponding accommodation coefficient was identified. The data on the orbital
decay of the ten satellites considered were taken by the authors from the historical
two-line elements (which can be found in Ref. [139]) provided by the United States
Space Surveillance Network. The meaning of the two-line elements is illustrated in
the following figure, which is due to the courtesy of NASA [140].
In orbit fitting, the ballistic coefficient (B = CDA/m) of each satellite was

adjusted to force the model of atmospheric density (the Jacchia-Bowman 2006
model was used) to agree with the air drag revealed by the tracking data resulting
from the historical two-line elements. For each of the satellites considered, the
area-to-mass ratio, A/m, was assumed to be constant, and consequently the changes
in the ballistic coefficient were only due to changes in the drag coefficient.
In this regard, Pardini et al. observe that eight of the ten satellites considered in
the study were spherical smooth objects, whereas the remaining two, namely the
Clementine and the Student Nitric Oxide Explorer (SNOE), were attitude-stabilised,
so as to have constant orientation with respect to the airstream.
Saunders et al. [141] observe that the drag term value (B*) given in the two-line
element set for a specific satellite is not the true ballistic coefficient defined above
(B = CDA/m, expressed in m2/kg) of the satellite itself. The constant conversion
factor given by Vallado [142] is
B ¼ 12:741621B
where B* is the drag term value taken directly from the two-line elements.
Vallado, in turn, took the value of the conversion factor from Hoots and
Roehrich [143], who gave this value in the following form
2
2:461 105 6378:135
The major semi-axis root mean square residuals, R, were computed by Pardini
et al. as follows
2 2 312
PN observed computed
a a
6 i¼1 i i 7
R¼4 5
N
where aobserved
i and acomputed
i are, respectively, the observed and the computed major
semi-axis of the satellite at the same epoch, and N is the number of observations
available, that is, the number of two-line elements used in the fitting. The complete
set of results found by Pardini et al. (for each of the ten satellites considered in the
study, a table of observed CD values is given as a function of the accommodation
coefficient a) can be found in Ref. [138]. A shortest account is given below. The
values of the observed drag coefficient range from 2.08 (relating to the Cosmos
1179, Cosmos 1427, and Cosmos 1616 satellites, at the altitude of 200 km) to 2.57
(relating to the Clementine satellite, at the altitude of 630 km). The corresponding
values of the accommodation coefficient range from 1.00 to 0.875. The conclusions
reached in the study are the following:
• energy accommodation coefficients are higher at sunspot maximum than at
minimum, and decrease more slowly with increasing altitude, which result is
consistent with the increased amount of oxygen that would be adsorbed on
satellite surfaces at solar maximum;
• consequently, drag coefficients are generally lower at solar maximum than at
minimum; and
• when the spherical shape of satellites is modified by attaching solar panels or
other objects, the drag coefficient increases substantially.
The following (and last) part of the present section is concerned with the effects
of atmospheric drag on satellite orbits.
The principal effect of atmospheric drag on an artificial satellite orbiting around
the Earth is the decrease of the mechanical energy possessed by the satellite. Since
atmospheric drag is greatest at perigee, where the satellite velocity and the atmo-
spheric density reach their maximum values in the course of each revolution, then
the loss of mechanical energy is also greatest at this point. Under the effect of this
braking impulse at perigee, an elliptic orbit reduces the height of its apogee in each
revolution, as shown in the following figure, which is due to the courtesy of NASA
[144], and the satellites spirals inward, until what is left of it eventually crashes on
the surface of the Earth.
For an elliptical orbit of high eccentricity, the perigee radius remains nearly
constant, whereas the apogee radius decreases with time, causing the major
semi-axis and eccentricity of the orbit to decrease. Consequently, the major
semi-axis, a, and the eccentricity, e, are of all orbital elements the most affected by
the in-plane drag force.
The variations of a and e caused by atmospheric drag are secular effects leading
to orbit shrinking and circularisation. In addition, Blitzer [1] notes that the cross
force due to atmospheric rotation causes small changes in i, X, and x.
In order to maintain a desired orbit, an artificial satellite subject to atmospheric
drag must perform altitude correction manoeuvres at regular intervals of time. Such
manoeuvres imply fuel consumption, as will be shown in the next chapter.
3.20 The Lifetime of an Earth Satellite Subject

to Atmospheric Drag
With reference to the following figure, we consider here the effects which the
atmospheric drag induces on the orbit of an Earth satellite.
Of course, these effects can be evaluated by means of the Lagrange planetary

equations, as the sequel will show. For the moment, in accordance with Meier
[145], we consider an artificial satellite, whose mass m is much smaller than the
mass M of the Earth, which revolves about the Earth in an elliptical orbit of major
semi-axis a.
Let v, fD, and fG be, respectively, the velocity vector of the satellite, the air-drag
force, and the gravity force acting on the satellite at a given time t. According to the
second principle of dynamics, the resulting force acting on the satellite is
GMm r lm r
f ¼ ma ¼ f G þ f D ¼ 2
þ fD ¼ 2 þ fD
r r r r
where G is the gravitational constant, l = GM is the gravitational parameter of the

Earth, and r is the geocentric position vector of the satellite.
The magnitude of the orbital velocity results from the vis-viva integral (Chap. 1,
Sect. 1.1), that is,

1
2 1 2
v¼ l
r a
and the magnitude of the atmospheric drag force is

1 2
fD ¼ CD qv A
2
where the atmosphere has been assumed to rotate at the same speed as that of the
solid Earth. This assumption is an approximation to the truth, as will be shown in
the next section. The total mechanical energy, E, of the satellite is the sum of its
kinetic and potential energies, that is
3.20 The Lifetime of an Earth Satellite Subject to Atmospheric Drag 499

1 2 lm 1 2 1 lm lm
E ¼ mv ¼ m l ¼
2 r 2 r a r 2a
The orbital period of the satellite is

3 12
a
T ¼ 2p
l
The elementary work, dL, done by the drag force in a time interval dt is
dL ¼ f D ds ¼ dE
where ds is the elementary displacement of the satellite in the time interval dt.
The work, dL, done by the drag force in the time interval dt is equal to the loss,
dE, of mechanical energy of the satellite in the same interval.
The loss of mechanical energy in the time interval dt is
dE d½lm=ð2aÞ lm da
¼ ¼ 2 ¼ fD v
dt dt 2a dt
Solving the preceding expression for da/dt yields
da fD v ½CD ð1=2 qv2 ÞAv a2 CD A 3 a2

¼ ¼ ¼ qv ¼ Bqv3
dt lm=ð2a2 Þ lm=ð2a2 Þ l m l
The preceding equation is usually written da/dt = –(a2/l)Bqv3, because the

satellite major semi-axis decreases with time as a result of the atmospheric drag.
The rate of change of the orbital period T of the satellite in the time interval dt
results from differentiating T with respect to time, as follows
h 1
i
dT d 2pða =lÞ
3 2
2p 3 3 da 3 1 da
¼ ¼ 1 að21Þ ¼ T
dt dt l2 2 dt 2 a dt
Substituting da/dt = (a2/l)Bqv3 into the preceding expression leads to
dT 3 1 da 3 1 a2 3a
¼ dt ¼ B q v3 dt ¼ B q v3 dt
T 2 a dt 2a l 2l
Hence, the relative change in orbital period over one revolution is
ZT
DT 3a
¼ B qv3 dt
T 2l
0
Therefore, DT/T depends on the ballistic coefficient B = CDA/m. When the

orbital elements and the ballistic coefficient of a satellite are known, then the
average atmospheric density can be determined.
We have shown above the effects of atmospheric drag on the major semi-axis
and period of the perturbed orbit. Let us consider now the effects on eccentricity.
If we assume that the perigee radius, rP, remains nearly constant, then
drP d½að1 eÞ

¼ ¼0
dt dt
Hence
da de
ð 1 eÞ a ¼0
dt dt
By solving the preceding expression for de/dt, there results
de 1 e da
¼
dt a dt
where da/dt has been determined above. Therefore, a decrease in major semi-axis
with time due to atmospheric drag also causes a decrease in eccentricity.
First-approximation values of major semi-axis and eccentricity can be computed
by using Taylor-series expansions truncated after the first-order terms, as follows

2
da a
a a0 þ Dt ¼ a0 0 Bq0 v30 Dt
dt 0 l

de 1 e0
e e0 þ Dt ¼ e0 a0 Bq0 v30 Dt
dt 0 l
The expressions indicated above form the basis for computing the lifetime of a
satellite revolving about the Earth, in the absence of trim burns.
A review of several methods meant to this purpose has been performed by de
Lafontaine and Garg [146]. A brief account of this review is given below.
Since the equations of motion for an Earth satellite subject to atmospheric drag
and gravitational perturbations cannot be solved analytically, then simplifying
methods are necessary to obtain approximate analytical solutions. The most com-
mon of such methods are:
(a) neglecting terms of small order;
(b) orbital averaging;
(c) asymptotic expansions;
(d) perturbation solutions; and
(e) variation of parameters.
The first of these methods is the most widely used. For example, in case of orbits
of small eccentricity (e 1), powers en with n greater than or equal to two are
usually neglected. For such near-circular orbits, the vis-viva integral

2 1
v2 ¼ l
r a
is approximated to v2 = l/r. In addition, the powers f n

of the Earth flattening
aE bE
f ¼
aE
(where aE and bE are the semi-axes of the ellipsoid approximating the Earth) are
also neglected when n is greater than or equal to two.
The orbital averaging method is based on neglecting the variations of the orbital
elements a, e, i, X, and x over an orbital period. By so doing, mean values of the
orbital elements named above are taken within each revolution, because their
changes in that time are very small in comparison with the change of orbital period.
Thus, the orbital motion comprises a fast-varying term (the true anomaly) and a set
of slow-varying terms (the orbital elements a, e, i, X, and x).
The mathematical theory used for this purpose is the Krylov–Bogoliubov
method, whose final objective is to eliminate the short-term effects of perturbations
and derive the elements affected only by the long period and the secular pertur-
bations. This theory was originally applied to the analysis of oscillating processes in
nonlinear mechanics. Subsequently, applications of this method to the orbits of
artificial satellites have been made by Struble [147], Kyner [148], and Musen [149].
The method of asymptotic expansions, also called generalised method of aver-
ages or Krylov–Bogoliubov–Mitropolsky method, is an extension of the method
described above. According to Cefola [150], the conventional equations of motion
are replaced with those written for the mean equinoctial elements (see Chap. 1,
Sect. 1.10) and the short periodic expressions. These expressions are Fourier series
with slowly varying coefficients and trigonometric variables related to the satellite
phase angle (which is the angle between the Sun and the satellite as seen by the
observer on the surface of the Earth) and the rotation of the Earth. The slowly
varying coefficients are evaluated at the output times by using low-order interpo-
lators. The mean elements can be estimated directly from the tracking data, by using
appropriate filtering techniques [150]. Among various authors who have applied
this method to the computation of orbits of Earth satellites subject to atmospheric
drag, we cite particularly Zee [151] and Barry et al. [152].
The method of perturbation solutions linearises the equations of motion around a
reference solution and provides analytical solutions of the linearised equations. This
method has been applied by Perkins [153].
The method of variation of parameters, due to Euler and Lagrange, computes the
change of the orbital elements due to the perturbing forces acting on the satellite.
This method will be considered at the end of the present section.
According to de Lafontaine and Garg [146], the lifetime of an Earth satellite

subject to atmospheric drag can be divided into two phases. They are:
• decay stage; and
• re-entry stage.
The decay stage goes from the moment of satellite injection into its initial orbit
to the moment in which the mechanical energy possessed by the satellite becomes
so low as to preclude the completion of an orbit around the Earth, at which time
re-entry begins. During this stage, the perigee altitude decreases from its initial
maximum value to a value going from 120 to 150 km. This stage is by far the
greater portion of the total lifetime. The rate of variation of the orbital elements a, e,
i, X, and x of the satellite is very small at the beginning of this stage and becomes
by degrees less small as the satellite approaches the end of the stage.
The re-entry stage goes from the end of the decay stage to the end of the satellite
life. The duration of this stage ranges from a few minutes to a few hours. The
perigee altitude decreases from the critical value mentioned above (120–150 km) to
a value at which the satellite is destroyed through volatilisation due to heat or
impact with the surface of the Earth. During this stage, the satellite follows not an
orbit but a ballistic trajectory characterised by a large rate of variation of the orbital
elements a, e, i, X, and x. Unless re-entry is caused by an impulsive manoeuvre
(such as the firing of speed-reducing thrusters) performed to that effect, this phe-
nomenon does not occur instantaneously, but rather gradually. Consequently, it is
difficult to determine a critical atmospheric density or a critical perigee altitude at
which re-entry takes place.
However, a fixed perigee height or a fixed density value has been chosen by
various authors to mark the end of life for a satellite.
For example, in case of an artificial satellite revolving in a circular orbit around
the Earth, its lifetime is defined by Garcia [144] as the time that it takes the satellite
to decay from its initial altitude to an altitude corresponding to a density of about
5 10−8 kg/m3. In other words, re-entry is assumed to occur when the satellite has
descended from its initial altitude to the altitude corresponding to the value
q = 5 10−8 kg/m3 of atmospheric density. According to the US Standard
Atmosphere, 1976 [154], this altitude is about 114 km.
King-Hele [155, 156] defines a critical radius of perigee, r*, measured from the
centre of the Earth to the satellite, equal to 6510 km, such that r* is the radius of the
last orbit of the satellite. In case of orbits of small inclination with respect to the
equatorial plane, r* corresponds to an altitude hmin = 132 km. This value is used in
the present book.
The method shown below to compute the lifetime of an Earth satellite subject to
drag is due to Garcia [144]. We consider separately the two cases of circular (or
nearly circular) orbit and elliptical orbit.
In case of an artificial satellite revolving around the Earth in a circular orbit, the
change in mechanical energy due to drag is
I Z2p
1
DE ¼ fD d‘ ¼ CD q A v2c rdh
2
0
where

2 1 2 1 l
v2c ¼l ¼l ¼
r a r r r
is the orbital velocity of the satellite, and the closed-path integral

I
fD d‘
is taken along a circle of radius r. The atmosphere has been assumed here to rotate
at the same angular velocity as the solid Earth. The quantity h is the angle swept out
by the radius vector, r, with respect to a direction of reference, as shown in the
following figure.
Assuming a constant attitude of the satellite with respect to the air flow, the
projected area, A, of the satellite normal to the velocity vector is also constant in
time. Consequently, the preceding expression reduces to
Z2p
1
DE ¼ lCD A qdh
2
0
The atmospheric density, q, is assumed to vary with the geometric altitude, h,

above the Earth as follows

h h0
q ¼ q0 exp
H
where q0 is the atmospheric density at a height of reference, h0, and H (which is such
that h/H be much greater than unity) is the density scale height, defined as follows
q
H¼
dq
dh
[146]. Under this hypothesis, the change in mechanical energy due to drag becomes
Z2p

1 h
DE ¼ lq0 CD A exp dh
2 H
0
In addition, assuming a small change in altitude in each revolution, the change in

energy is

h
DE ¼ lq0 CD A p exp
H
The total mechanical energy, E, of a satellite of mass m is the sum of its kinetic
and potential energies, that is

1 2 lm 1 2 1 lm lm
E ¼ mv ¼ m l ¼
2 r 2 r r r 2r
Consequently, the change in the total mechanical energy, for a constant mass of
the satellite, is
l ml
DE ¼ mD ¼ Dr
2r 2r 2
Since the change in energy is equal to the work done by the drag force, then
introducing DE = –l q0 CD A p exp(–h/H) into the expression written above leads to
ml

h
Dr ¼ lq0 CD Ap exp
2r 2 H
which, solved for Dr, leads to

CD A h h
Dr ¼ 2pr 2 q0 exp ¼ 2pr 2 q0 B exp
m H H
The preceding equation expresses the change in radius vector in one revolution
(N = 1) of the satellite about the Earth. Since the value of Dr is small, then the first
variation may be taken as the derivative, and the preceding expression may be
rewritten as follows

h
ðRE þ hÞ2 exp dh ¼ 2pBq0 dn
H
where RE is the mean radius of the Earth, r = RE + h, N is the number of revolu-

tions performed by the satellite, and n is a variable used in the integration of
N. Now, (RE + h)−2 = R−2
E (1 + h/RE)
−2
is expanded in a Maclaurin series of h/RE,
as follows
1 00 1 h ð3Þ i 3
f ð xÞ ¼ ½f ð xÞ0 þ ½f 0 ð xÞ0 x þ ½f ð xÞ0 x2 þ f ð xÞ x þ . . .
2! 3! 0
where, for convenience, h/RE has been set equal to x.

The first, second, and third derivatives of f(x) = (1 + x)−2 are
f 0 ðxÞ ¼ 2ð1 þ xÞ3

f 00 ðxÞ ¼ 6ð1 þ xÞ4
f 000 ðxÞ ¼ 24ð1 þ xÞ5
The function f(x) and its derivatives, evaluated at x = 0, are respectively
½f ðxÞ0 ¼ 1
½f 0 ðxÞ0 ¼ 2
½f 00 ðxÞ0 ¼ 6
½f 000 ðxÞ0 ¼ 24
Therefore
6 2 24 3
ð1 þ xÞ2 ¼ 1 2x þ x x þ ¼ 1 2x þ 3x2 4x3 þ
2! 3!
Since x = h/RE, then

"

2
3 #
2 h 2 h h h
ðRE þ hÞ ¼ R2
E 1þ 2
¼ RE 1 2 þ3 4 þ
RE RE RE RE
Now we truncate the expansion of h/RE after the third power and substitute
"

2
3 #
2 h h h
ðRE þ hÞ R2
E 12 þ3 4
RE RE RE
into

2h
ðRE þ hÞ exp dh ¼ 2pBq0 dn
H
This yields
"

2
3 #

h h h h
R2
E 12 þ3 4 exp dh ¼ 2pBq0 dn
RE RE RE H
After integrating the left-hand side of the preceding equation from h2 to h1 and
the right-hand side from zero to N (where N is the number of revolutions performed
by the satellite to decay from h1 to h2), there results
H
R2E

¼ 2pBq0 N
h1 h2
k1 exp k2 exp
H H
where

2

2 6H 24H 2 h1 12H h1 3
k1 ¼ 1 ðh1 H Þ þ 2 þ 3 þ 3þ 4
RE RE RE RE RE RE

2

3
2 6H 24H 2 h2 12H h2
k2 ¼ 1 ðh2 H Þ þ 2 þ 3 þ 3þ 4
RE RE RE RE RE RE
It is necessary now to obtain an average orbital period, ⟨T⟩, to get an expression

for lifetime in terms of time, as follows
N hT i ¼ t
By using the first mean-value theorem for integrals, Garcia obtains
Zr1
2
2p 3 4 p r1 r2
hT i ¼ r 2 dr ¼ 2
l2
1
5 r1 r2 vc1 vc2
r2
where vc1 and vc2 are the velocities of a satellite in circular orbits of radii,
respectively, r1 and r2, that is, vc1 = (l/r1)½ and vc2 = (l/r2)½.
The lifetime expression is obtained as follows. The time necessary to atmo-

spheric drag to reduce the initial eccentricity of the orbit to zero results from
equation N⟨T⟩ = t, where

2
4 p r1 r2
hT i ¼ 2
5 r1 r2 vc1 vc2
and N is given by the expression derived above, which is rewritten below for
convenience
H
R2E

¼ 2pBq0 N
h1 h2
k1 exp k2 exp
H H
This leads to the following expression of the satellite lifetime

2

2 r1 r2 H h1 h2
t¼ 2 k 1 exp k 2 exp
5Bðh1 h2 Þ vc1 vc2 q0 RE
2 H H
In the computation, q0 and H are assumed to be constant between h1 and h2.

Therefore, the interval [h1, h2] must be chosen so as to satisfy this condition. If
necessary, the interval [h1, h2] must be split up into sub-intervals, so as to satisfy
piecewise the condition for each of them. The preceding equations express the
satellite lifetime, t, in seconds. To have t in terms of years, it is necessary to divide
the lifetime computed above by the number of seconds contained in a year.
In case of a computation performed by hand, Garcia’s formulae require much
labour. Instead of them, it is possible to use those given below.
The lifetime, t, of an Earth satellite in a circular orbit is expressed by King-Hele
[132, p. 62, formula 4.89] as follows
gHT
t¼
2pBqr 2
where

r r
g ¼ 1 exp
H
T is the orbital period, H is the density scale height corresponding to the initial
geometric altitude of the circular orbit, B is the ballistic coefficient of the satellite, q
is the atmospheric density corresponding to the same altitude, r is the initial geo-
centric radius of the orbit, and r* = 6510 km is the critical geocentric radius. An
example of application of these formulae will be given below.
In case of an artificial satellite revolving around the Earth in an initially elliptical
orbit, the method described by Garcia proceeds as follows.
First, the initial (t = t0) orbit is said to be of appreciable eccentricity (e0), when
the following condition is satisfied at t = t0
x0
[2
H
where H is the density scale height computed at the altitude of the initial perigee,
x0 = ½(hA0 − hP0), and hA0 and hP0 are the altitudes of the satellite at, respectively,
initial apogee and initial perigee. In this case, the satellite lifetime is
vcA0 H x0
t¼ f ; e0
lE BqP0 H
where
x0 x
0
x exp 11 29 2
f
0
; e0 ¼ H H
x exp ðe0 Þ e0 exp ðe0 Þ þ e exp ðe0 Þ
H 2I 1
0 6 16 0
H
7He0 expðe0 Þ
þ
8x0
In the preceding expression, only the terms of the order of magnitude of e20 have
been retained, whereas those containing powers of e0 with exponents greater than
two have been neglected. In addition, a0 is the major semi-axis of the initial
elliptical orbit, vcA0 = (lE/a0)½ is the velocity of a circular orbit of radius a0, and
Zp
1
I 1 ðzÞ ¼ ðcos hÞexpðz cos hÞdh
p
0
is the modified Bessel function of the first kind, of order n = 1.

If the argument of I1(z) is much greater than unity, then I1(z) can be computed by
using the following approximation [157, p. 377, formula 9.7.1]:
" #
expðzÞ 4 1 ð4 1Þð4 9Þ ð4 1Þð4 9Þð4 25Þ
I 1 ðzÞ 1 1 þ þ
ð2pzÞ2 8z 2!ð8zÞ2 3!ð8zÞ3
The preceding equations express the satellite lifetime, t, in seconds. To have t in

terms of years, it is necessary to divide the lifetime computed above by the number
of seconds contained in a year. In summary, the computation of the lifetime of an
artificial satellite, revolving about the Earth along an initially elliptical orbit and
subject to the perturbing action of aerodynamic drag, may be performed in two
steps. The first step is the computation of the time te taken by drag to decrease the
height of the satellite apogee, so as to place the satellite into a nearly circular orbit.
The second step is the computation of the time tc spent by the satellite in a decay
spiral going from the altitude of the nearly circular orbit determined in the first step
to the minimum altitude hmin = 132 km (corresponding to the radius
r* = 6510 km). The total lifetime t of the satellite results from the sum te + tc. As
an example of application of the theory shown above to a practical case, let us
consider the Cosmos 1179 (also known as Taifun 1-B) satellite. According to
NASA [158], the data concerning this satellite and its initial orbit are the following:
Shape: sphericalðd ¼ 2 mÞ Launch date: 14 May 1980 12 : 57: 00 UTC

Mass: m ¼ 650:0 kg Decay date: 17 July 1989
Initial perigee: hP0 ¼ 310:0 km Initial orbit inclination: i ¼ 83 :0
Initial apogee: hA0 ¼ 1570:0 km Initial eccentricity: e0 ¼ 0:08603
Initial period T0 ¼ 103:5 min ¼ 6210 s
The planetary constants used here are lE = 3.986 1014 m3/s2 and
RE = 6.371 106 m. The drag coefficient used in the computation of the whole
lifetime of this satellite is CD = 2.2, as recommended by King-Hele and Walker
[156] for perigee altitudes in the range 150–400 km.
The projected area of the satellite normal to the relative velocity vector is
pd 2 3:1416 22
A¼ ¼ ¼ 3:1416 m2
4 4
The area-to-mass ratio of the satellite is
A 3:1416
¼ ¼ 0:0048332 m2 =kg
m 650
The major semi-axis of the initial orbit results from Kepler’s third law

13
1
lE T02 3:986 1014 62102 3
a0 ¼ ¼ ¼ 7:3022 106 m
4p2 4 3:14162
The average altitude of the satellite in its initial orbit is
hA0 hP0 1570000 310000

x0 ¼ ¼ ¼ 630000 m
2 2
By interpolating the values given in the tables of the US Standard Atmosphere,

1976 [154], the pressure scale height corresponding to the initial altitude of perigee
(hP0 = 310000 m) results
HP0 ¼ 53741:6 m
This value refers to the pressure scale height. However, as has been shown by
Prölss [159], at altitudes greater than 200 km, the difference between the density
scale height and the pressure scale height is negligible.
The atmospheric density at perigee results
qP0 ¼ 1:552 1011 kg=m3

Consequently, the argument of the modified Bessel function of the first kind, of
order n = 1, is
x0 630000
¼ ¼ 11:723
HP0 53741:6
The corresponding modified Bessel function of the first kind, first order, results
from the following series expansion
"
expð11:723Þ 41 ð4 1Þð4 9Þ
I 1 ð11:723Þ 1 þ
ð2 3:1416 11:723Þ
1
2 8 11:723 2!ð8 11:723Þ2
#
ð4 1Þð4 9Þð4 25Þ
¼ 13899:3
3!ð8 11:723Þ3
The velocity of a satellite in a circular orbit of radius a0 is

12
12
lE 3:986 1014
vcA0 ¼ ¼ ¼ 0:73883 104 m/s
a0 7:3022 106
The quantity f(x0/HP0, e0) results from

x0 x0

exp
x0 HP0 HP0 11 29 2
f ; e0 ¼
exp ðe0 Þ e0 exp ðe0 Þ þ e exp ðe0 Þ
HP0 x0 6 16 0
2I1
HP0

7HP0 e0 exp ðe0 Þ
þ
8x0
After substituting x0/H = 11.723, I1(11.723) = 13899.3, and e0 = 0.08603 into

the preceding equation, we obtain

x0
f ; e0 ¼ 48:885
HP0
and the time te taken by atmospheric drag to change the elliptical orbit of the
satellite into a circular orbit, through the decrease of the apogee height, is
vcA0 HP0 x0 0:73883 104 53741:6

te ¼ f ; e0 ¼ 48:885
lE BqP0 H 3:986 1011 2:2 0:0048332 1:552 1011
¼ 2:9508 108 s ¼ 9:3508 years
The remaining lifetime (tc) of the satellite is computed as suggested by

King-Hele [132]:
gHT
t¼
2pBqr 2
where

r r
g ¼ 1 exp
H
Remembering that the orbital period of a satellite in a circular orbit of radius

r around the Earth is

3 12
r
T ¼ 2p
lE
the remaining lifetime may also be expressed as follows
gH
tc ¼ 1
BqðlE r Þ2

r r 6688000 6510000
g ¼ 1 exp ¼ 1 exp ¼ 0:96356
HP0 53741:6
Therefore, the remaining lifetime is given by
gH 0:96356 53741:6
tc ¼ 1 ¼ 1
BqðlE r Þ2
2:2 0:0048332 1:552 1011 ð3:986 1014 6:688 106 Þ2
¼ 6:0775 106 s ¼ 70:341 days ¼ 0:19249 years
The total lifetime of the Cosmos 1179, resulting from this computation, is then
t ¼ te þ tc ¼ 9:3508 þ 0:19249 ¼ 9:5433 years
By comparison, assuming the satellite to have been injected into its first orbit at
the launch time, and a time of decay equal to 12:00:00 UTC, the real lifetime is
2447725:0 2444374:039583
t¼ ¼ 9:1744 years
365:25
The methods shown above compute analytically the lifetime of an Earth satellite
subject to the perturbing action of atmospheric drag. According to Garcia, the
maximum deviation of the results obtained by using his method from those
obtained by means of numerical integration is less than 13% (or less than 25% in
case of integration of the Lagrangian planetary equations). These deviations only
occur at the higher altitudes and decrease rapidly with decreasing altitudes [144]. In
the example of computation given above, the relative error committed in computing
the satellite lifetime is
9:5433 9:1744
¼ 0:040204 4%
9:1744
The difference between the computed values and the real values depends mainly
on the following reasons:
• in the method shown above, the atmosphere and gravitational field of the Earth
are considered, both of them, perfectly spherical, whereas the atmospheric
oblateness and the odd zonal harmonics cause oscillations in the atmospheric
density, due to the variable altitude of perigee;
• semi-annual variations in atmospheric density, due to the variable Earth–Sun
distance in the course of an orbit of the Earth around the Sun;
• variations in atmospheric density caused by the eleven-year solar cycle (in times
of high solar activity, the satellite lifetime is shorter due to higher values of
atmospheric density);
• day-to-night variations in atmospheric density, due to the diurnal heating–
cooling cycle;
• rotation of the atmosphere with the Earth, due to the friction between the Earth
and its atmosphere (the ratio K = xatm/xE of the angular velocity of the
atmosphere to that of the Earth is neither equal to unity nor constant at any
altitude, as will be shown in the following sections); and
• other short-term density variations due to unpredictable changes in solar
activity.
In order to take account of the major effects, King-Hele and Walker propose to
compute a preliminary value of lifetime and to correct this value by means of
appropriate formulae. Particulars on this method can be found in Ref. [156].
The lifetime of an Earth satellite subject to atmospheric drag can also be com-
puted by numerical integration of the Lagrange planetary equations in Gaussian
form (see Sect. 3.2), which express the changes in the orbital elements. As shown
above, the life of a satellite comes to an end, when the major semi-axis of its orbit
reaches the critical value a* = 6510 km.
Such is the method used by Belcher et al. [160], which is shown below. To apply
this method to the case of interest, it is necessary to determine the three components
(radial, transverse, and bi-normal) of the perturbing acceleration due to atmospheric

drag, as follows
a ¼ ar ur þ a/ u/ þ ah uh
It is understood that the perturbing acceleration vector a considered below

fD 1 CD A
a aD ¼ ¼ q jvrel jvrel
m 2 m
is only the part of the perturbing acceleration which is due to atmospheric drag.
The three components (ar, a/, and ah) of a are then introduced into the Lagrange
planetary equations in Gaussian form
2a2 h p i
h r
1
h

r cosðx þ /Þ
i0 ¼ ah
h

0 r sinðx þ /Þ
X ¼ a/
h sin i

0 p cos / ðp þ r Þ sin / r sinðx þ /Þ cos i
x ¼ ar þ a/ ah
he he h sin i
b
M0 ¼ n þ ½p cos / 2er ar ½ðp þ r Þ sin /a/
ahe
where an apex denotes the time derivatives of the orbital elements.

As has been shown above, the relative velocity vector vrel, which appears in the
expression of the perturbing drag acceleration, is the velocity of the satellite with
respect to the atmosphere, that is,
vrel ¼ v vatm
where v is the velocity of the satellite with respect to the inertial reference system,
and vatm is the velocity vector of the Earth atmosphere, with respect to the same
reference system, at the satellite position. The magnitudes of v and vatm are

1
2 1 2
v¼ l vatm ¼ xE r sin h
r a
where xE is the angular velocity of both the Earth and its atmosphere, r is the radius
vector from the geocentre to the satellite, and h is the polar distance (or co-latitude) of
the satellite, as shown in the following figure, due to the courtesy of NASA [161].
The components (ar, a/, and ah) of the acceleration vector a due to air drag
1 CD A
a¼ q jvrel j2 urel
2 m
which has the direction of the unit vector urel, can be determined by projecting
a onto the directions of, respectively, ur, u/, and uh.
The inertial velocity vector v of the satellite can be expressed as follows
v ¼ vr ur þ v/ u/ ¼ ðv sin cÞur þ ðv cos cÞu/
where c is the elevation angle, shown in the following figure, of the inertial velocity
vector v above a straight line perpendicular to the satellite position vector r and
contained in the orbital plane.
The velocity vector of the atmosphere is
vatm ¼ xE r ¼ xE uzE r ur ¼ xE ðsin u sin i ur þ cos u sin i uf þ cos i uh Þ r ur

¼ ðxE r cos iÞ u/ ðxE r cos u sin iÞ uh
where u = x + / is the argument of latitude at epoch (see Chap. 1, Sect. 1.9).

Hence, the relative velocity vector of the satellite is
vrel ¼ v vatm ¼ ðv sin cÞur þ ðv cos cÞu/ ðxE r cos iÞu/ þ ðxE r cos u sin iÞuh
¼ ðv sin cÞur þ ðv cos c xE r cos iÞu/ þ ðxE r cos u sin iÞuh
Consequently, the perturbing acceleration vector a due to air drag
fD 1 CD A
a aD ¼ ¼ q jvrel jvrel
m 2 m
may be written as follows
1 CD A
aD ¼ q jvrel j½ðv sin cÞur þ ðv cos c xE r cos iÞu/
2 m
þ ðxE r cos u sin iÞuh
and its three components along the radial, transverse, and bi-normal directions are
respectively
1 CD A
aDr ¼ q jvrel jv sin c
2 m
1 CD A
aD/ ¼ q jvrel jðv cos c xE r cos iÞ
2 m
1 CD A
aDh ¼ q jvrel jxE r cos u sin i
2 m
The components of the inertial velocity vector, v, of the satellite can be

expressed in terms of the orbital elements and the true anomaly, /, as follows

12
l
vr ¼ v sin c ¼ e sin /
p

12
l
v/ ¼ v cos c ¼ ð1 þ e cos /Þ
p
This is because
v ¼ r0 ¼ ðrur Þ0 ¼ r 0 ur þ r/0 u/ ¼ vr ur þ v/ u/
where the prime sign denotes time derivatives. Remembering that

2

p h 1
r¼ ¼
1 þ e cos / l 1 þ e cos /
the radial (vr) and transverse (v/) components of the velocity vector v are
"
1 # "
1 #
0l 2 l 2
vr ¼ e sin / ur þ ð1 þ e cos /Þ u/
p p
Therefore, the three components (along, respectively, the radial, transverse, and
bi-normal directions) of the drag acceleration vector are expressible as follows

12
1 CD A l
aDr ¼ q jvrel j e sin /
2 m p
"
1 #
1 CD A l 2
aD/ ¼ q jvrel j ð1 þ e cos /Þ xE r cos i
2 m p
1 CD A
aDh ¼ q jvrel jxE r cosðx þ /Þ sin i
2 m
where
h i12
jvrel j ¼ ðv sin cÞ2 þ ðv cos c xE r cos iÞ2 þ ðxE r cos u sin iÞ2
Taking account of the expressions of v sin c and v cos c shown above, the
preceding expression may be written as follows
8 "
1 #2 91
<l l 2 =2
2
jvrel j ¼ e sin / þ
2 2
ð1 þ e cos /Þ xE r cos i þ ½xE r cos u sin i
:p p ;
According to Belcher et al. [160, appendix B, p. 114], the preceding expression

can be approximated as follows

12
l 1 pxE cos i
jvrel j 1 þ e2 þ 2e cos / 2 1
p ð1 þ e2 þ 2e cos /Þ2
The three components of the drag acceleration vector are introduced into the
Lagrange planetary equations in Gaussian form, and these equations can be inte-
grated numerically.
The satellite altitude h(/) above the terrestrial ellipsoid (which approximates the
geoid) is expressible as the difference between the radius vector r(/) of the satellite
and the radius rE(/) of the ellipsoid, as follows

p R2E 1 e2E
hð/Þ ¼ r ð/Þ rE ð/Þ ¼
1 þ e cos / 1 e2E 1 sin2 i sin2 ðx þ /Þ
where RE and eE are, respectively, the major semi-axis and the eccentricity of the
terrestrial ellipsoid.
3.21 The Fundamental Properties of the Earth

Atmosphere
By atmosphere we mean the gaseous envelope surrounding a celestial body. In the

case of the Earth, the greatest part of this envelope is made up of nitrogen (78% by
volume) and oxygen (21% by volume). The remaining part (1%) of the atmosphere
is made up of the so-called trace gases, because they are present in very small
concentrations. The most abundant of them is argon (0.93%), which, as is the case
with all other noble gases (neon, helium, krypton, and xenon), is inert and does not
react chemically with the other atmospheric gases.
Hydrogen is also one of the trace gases, and is contained in the atmosphere to a
negligible extent (0.000056%) because of its very low molecular weight, which
makes it particularly apt to escape the gravitational force of the Earth. This mixture
of gases is called air. The following table (from Schlatter [131]) gives the con-
centration of gases constituting dry air near sea level in parts per million by volume
(ppmv), in order of decreasing concentration. The mean molecular weight is
expressed in kg/kmol.
Gas name Chemical symbol Mean molecular weight Concentration

Nitrogen N2 28.013 780840
Oxygen O2 31.999 209460
Argon Ar 39.948 9340
Carbon dioxide CO2 44.010 384
Neon Ne 20.180 18.18
Helium He 4.003 5.24
Methane CH4 16.043 1.774
Krypton Kr 83.798 1.14
Hydrogen H2 2.016 0.56
(continued)
(continued)
Gas name Chemical symbol Mean molecular weight Concentration
Nitrous oxide N2O 44.012 0.320
Xenon Xe 131.293 0.09
Ozone O3 47.998 0.01–0.10
The Earth atmosphere is made up of several concentric layers, also called strata.
Each of them is separated from the contiguous ones by narrow zones and has its
own chemical composition and temperature distribution. These layers and their
temperatures are shown in the following figure (courtesy of COSPAR [162]).
More than 99% of the total mass of the atmosphere is confined in the altitude
range going from the Earth surface to 40 km.
The atmospheric stratum nearest to the surface of the Earth is the troposphere,
which contains the greatest amount of gaseous mass and 99% of the water vapour
existing in the atmosphere. The troposphere begins at the Earth surface and ends at
a height variable, as a function of latitude and season of the year, from 8 to 18 km
[163]. The maximum height refers to the equator in the summer, whereas the
minimum refers to the poles in the winter. The density of the gases in the tropo-
sphere and the temperature decrease, both of them, with height. The temperature
3.21 The Fundamental Properties of the Earth Atmosphere 519
range goes from a mean value of 17 °C at the bottom to a mean value of −51 °C at
the top of the troposphere [164]. Almost all weather phenomena occur in this
region. The transition zone from the troposphere to the next layer, the stratosphere,
is called the tropopause. In this transition zone, temperature remains constant with
increasing altitude.
The stratosphere extends from 10 to 50 km above the Earth surface. Here the
temperature remains nearly constant up to an altitude of 55 km, and then increases
gradually to 200–220 K at the lower level of the stratopause (50 km), at which
altitude it begins to decrease.
The ozone layer is at an altitude between 20 and 30 km. Here, the concentration
of ozone is about 10 parts per million by volume, which value is to be compared
with 0.04 parts per million by volume in the troposphere.
The following figure, due to the courtesy of the US Government, illustrates the
temperature-altitude profile for the U.S. Standard atmosphere, 1976 [154] in the
range 0–100 km of geometrical altitude.
This figure shows the increase of temperature in the altitude range of the
stratosphere (10–50 km). At the altitude of about 50 km, which marks the begin-
ning of the stratopause, the temperature remains constant (about 0 °C or 273 K) in
the altitude range 50–55 km, that is, within the stratopause. Above the stratopause,
the atmospheric temperature decreases up to a value of about 180–190 K.
The next layer above the stratopause is the mesosphere, which extends up to an
altitude of about 90 km. As has been shown above, the mesosphere is characterised
by a decreasing temperature. Here, the concentrations of ozone and water vapour
are negligible. The mesosphere terminates with a narrow zone of separation, called
mesopause, from the following layer.
The altitude of about 90 km from the Earth surface marks the beginning of the
thermosphere, also known as the upper atmosphere, which extends up to an altitude
of about 600 km from the Earth surface. Here, the temperature generally increases
with altitude up to 1000–1500 K. According to Doornbos [165], the thermosphere
is characterised by the absorption of solar extreme-ultraviolet radiation, which

causes temperature (T) to increase asymptotically with altitude (h). In other words,
in the range of altitudes going from 200 km to 600 km, the temperature gradient
(dT/dh) is so small as to be negligible (T T∞).
In the lower thermosphere, that is, in the range of altitudes going from 120 km to
200 km, temperature is assumed to vary with altitude according to the Bates (or
Bates-Walker) temperature profile
T ¼ T1 ðT1 T120 Þ exp½sðh h120 Þ
[159, p. 43], where T∞ is the exospheric temperature, T120 = 350 K is the tem-
perature at the altitude h120 = 120 km, and s = 0.021 km−1 is a quantity determined
empirically which decreases with T∞. By the way, the exospheric temperature is the
value which the exponential function representing T(h) above 120 km asymptoti-
cally approaches at heights above about 600 km, where the mean free path of the
particles exceeds the scale height. The US Standard Atmosphere, 1976 [154] takes
T∞ = 1000 K in mean solar conditions.
The atmospheric layer beyond the thermosphere (from about 600 km upwards)
is the exosphere, which extends up to about 10000 km, where it becomes practi-
cally the same thing as interplanetary space. In this layer, atoms and molecules are
free to escape into space. A part of the upper atmosphere, known as the ionosphere,
comprises portions of the mesosphere, thermosphere, and exosphere, and extends
from about 50 km to more than 1000 km altitude. It contains particles (electrons
and electrically charged atoms and molecules) which are ionised by ultraviolet
radiation coming from the Sun. It is, in turn, divided into layers and has the
important property of reflecting radio waves coming from the Earth. The following
table, taken from Schlatter [131], gives some properties of neutral atmospheric gas
as a function of geometrical altitude.
h g n ⟨v⟩ t L ⟨M⟩ H
0 9.807 2.547 1025 458.9 6.919 109 6.633 10−8 28.96 8434
20000 9.745 1.849 1024 398.0 4.354 108 9.139 10−7 28.96 6382
40000 9.684 8.308 1022 427.8 2.104 107 2.034 10−5 28.96 7421
70000 9.594 1.722 1021 400.6 4.084 105 9.810 10−4 28.96 6570
100000 9.505 1.189 1019 381.4 2.68 103 1.42 10−1 28.40 6009
150000 9.360 5.186 1016 746.5 2.3 101 3.3 101 24.10 23380
200000 9.218 7.182 1015 921.6 3.9 100 2.4 102 21.30 36183
300000 8.943 6.509 1014 1080 4.2 10−1 2.6 103 17.73 51193
400000 8.680 1.056 1014 1149 7.2 10−2 1.6 104 15.98 59678
600000 8.188 5.950 1012 1356 4.8 10−3 2.8 105 11.51 88244
800000 7.737 1.234 1012 1954 1.4 10−3 1.4 106 5.54 193860
1000000 7.322 5.442 1011 2319 7.5 10−4 3.1 106 3.94 288200
3.21 The Fundamental Properties of the Earth Atmosphere 521
These data are the same as those of the US Standard Atmosphere, 1976 [154] for
dry air. The quantities shown in the table are geometric altitude (h, expressed in
metres), acceleration due to gravity (g, in m/s2), number density (n, number of
molecules in a m3), mean molecular speed (〈v〉, in m/s), collision frequency (t, in s−1),
mean free path (L, in m), mean molecular weight (〈M〉, in kg/kmol), and scale height
(H, in m). The US Standard Atmosphere, 1976 [154] also gives density (kg/m3),
pressure (N/m2), density scale height (km), and pressure scale height (km) in graphical
form, as a function of geometric altitude (km), as shown below.
As shown in the following figure (from [154]), there are two atmospheric scale
heights: the density scale height (H, defined in Sect. 3.20) and the pressure scale
height Hp.
The pressure scale height is defined as follows

p
Hp ¼
dp
dh
The values of the latter are tabulated in the US Standard Atmosphere 1976 as a
function of geometric altitude. The relationship between the two scale heights is
1 1 1 dT
¼ þ
H Hp T dh
According to Prössl [159], usually (and certainly above an altitude of 200 km)
the second term of the sum on the right-hand side of the preceding expression is
small in comparison to the first, so that the values of the density scale height
correspond closely to those of the pressure scale height (H Hp).
3.22 Atmospheric Density Models
As has been shown in Sect. 3.19, the atmospheric drag acceleration acting on an
Earth satellite
fD 1 CD A
aD ¼ ¼ q jvrel jvrel
m 2 m
depends, inter alia, on the values of atmospheric density (q) in the points of the
satellite orbit. According to Marcos [166], air density in the thermosphere, that is, in
the altitude range 90–600 km, depends mainly on two causes. They are the solar
extreme-ultraviolet (EUV) radiation and the solar wind. The heating due to solar
EUV radiation is on average about 80% of the total amount of energy reaching the
thermosphere. The solar EUV radiation varies according to the 27-day solar rotation
period and also to day-to-day oscillations. It also varies according to geomagnetic
storms, but on larger temporal intervals.
There are of course other causes of density variation along a satellite orbit, such
as altitude, local solar time, and day of the year. The effect produced by heating is
an increase in atmospheric density, because heating generates expansion and this, in
turn, increases the air density existing at a fixed altitude. Consequently, the decay of
an artificial satellite having a given ballistic coefficient (CD A/m) is more rapid in
periods of maximum solar activity than is in periods of minimum.
At a fixed point of a satellite orbit, the atmospheric density depends essentially on
two environmental quantities: the solar proxy F10.7 and the geomagnetic activity index
ap. The solar activity is measured by the values of the solar proxy F10.7, which is the
daily radio flux emitted by the Sun at a frequency of 2800 MHz (corresponding to a
wavelength of 10.7 cm). The values of the solar proxy F10.7 are normally expressed in
solar flux units (1 SFU = 1 10−22 W m−2 Hz−1). According to Schatten [167],
these values can vary from about 50 (at solar minima) to about 300 (at solar maxima).
According to Doornbos [165], they can vary from below 70 to around 370.
The magnetic field disturbances induced by changes in the solar wind are
measured by the ap (or Kp, which is essentially the logarithm of ap) index, where the
subscript p stands for planetary.
The general relationship between ap and Kp values is given in Table 1 of [168],
and also in the following table, which is due to the courtesy of NOAA [169].
Kp = 0o 0+ 1− 1o 1+ 2− 2o 2+ 3− 3o 3+ 4− 4o 4+
ap = 0 2 3 4 5 6 7 9 12 15 18 22 27 32
Kp = 5− 5o 5+ 6− 6o 6+ 7− 7o 7+ 8− 8o 8+ 9− 9o
ap = 39 48 56 67 80 94 111 132 154 179 207 236 300 400
3.22 Atmospheric Density Models 523
The values of ap range from 0 (minimum geomagnetic activity) to 400 (maxi-

mum geomagnetic activity), and the corresponding values of Kp range from 0 to 9o.
The scale of Kp from 0 to 9 is expressed in thirds of a unit, so that 0o stands for 0.0,
0+ stands for 0.33, 1− stands for 0.67, 1o stands for 1.0, etc. The geomagnetic
index ap or Kp measures solar particle radiation by its magnetic effects [169]. To an
increase in each of the quantities F10.7 and ap, there corresponds an increase in the
atmospheric density at altitudes greater than about 120 km.
For the purpose of computing drag forces and lifetimes for Earth satellites, there
are numerous atmospheric density models, which provide estimates of the statistical
mean temperature, total mass density, and number density of each atmospheric
constituent as a function of position, local time, Universal Time, and solar and
geomagnetic indices [170].
Vallado and Finkleman [171] have performed a review of these models, and in
particular of those listed below.
1. U.S. Standard Atmosphere, 1976;
2. Variations of the Jacchia-Roberts models (J71, J77 and GRAM99);
3. COSPAR International Reference Atmosphere (CIRA90);
4. Mass Spectrometer-Incoherent Scatter (NRLMSISE-00);
5. Drag Temperature Models (DTM-94 and DTM-03);
6. Marshall Engineering Thermosphere (MET-88 and MET-99);
7. GOST Russian Models (GOST-84 and GOST-04); and
8. General Circulation Models (TIGCM and TIEGCM).
According to Vallado and Finkleman, these models are based on the following
assumptions:
• the amount of the solar extreme-ultraviolet (EUV) radiation is adequately
measured by the solar proxy F10.7;
• the lower atmosphere rotates with the Earth;
• the atmospheric composition at sea level is substantially the same as that
resulting from the table given in the preceding section;
• the change in mean molecular mass in the mixing region below 105 km is
caused by oxygen dissociation;
• the shape of the temperature profiles remains constant with respect to basic
models;
• the drag coefficient results from CD = 2 m|aD|/(qv2relA);
• the atmospheric model accounts for the variations in atmospheric drag due to the
diurnal variations, the tilt of the Earth, the 27-day solar rotation, the 11-year
solar cycle, the semi-annual and seasonal variations, the oscillations in the
11-year cycle, the upper atmospheric winds, the magnetic storms, the irregular
short periodic variations, the tides, the spatial variations of components of the
atmosphere, the appropriate lags for all factors affecting density, and all other
factors;
• the exact method for employing the existing atmospheric indices is known and
used;
• all the input parameters (such as solar flux, geomagnetic indices, etc.) are
available with known precision in the future;
• any interpolation of the indices introduces no errors into the solution;
• the proper times and related interpolations are used with each index; and
• the atmospheric model is appropriate for the orbital class of the satellite under
consideration.
The US Standard Atmosphere 1976 (described in Sect. 3.21) has been developed
jointly by NOAA, NASA, and the United States Air Force. It is a static repre-
sentation of the terrestrial atmosphere from the ground level to 1000 km, in a period
of moderate (i.e. with F10.7 in the range 75–150 SFU) solar activity. In the altitude
range going from the sea level to 71 km, this model assumes a mean molecular
weight and divides the atmosphere into six layers, as shown in the following table
(courtesy of the University of California [172]).
Layer No. Name Lower altitude (km) Upper altitude (km)

1 Troposphere 0 11
2 Stratosphere 11 20
3 – 20 32
4 – 32 47
5 – 47 51
6 Mesosphere 51 71
According to the model, at sea level the temperature is 288.15 K (15 °C), the
density is 1.225 kg/m3, and the pressure is 101325 N/m2.
Each of the six layers is defined by its own linear gradient of temperature. For
example, the troposphere has a temperature gradient of −6.5 K/km [131].
The US Standard Atmosphere 1976 assumes a constant value for the accelera-
tion of gravity (g). Consequently, it distinguishes between the ordinary geometric
altitude (h) and the geopotential altitude (z), defined as follows
Zh
1
z¼ gðu; zÞdz
g0
0
where u is the latitude and g0 is the acceleration of gravity at mean sea level.
The geometric altitude (h) results from the geopotential altitude (z) as follows
zRE
h¼
RE z
[154], where RE is the effective radius of the Earth at a specific latitude (the US
Standard Atmosphere 1976 takes RE = 6356766 m). This is because the gravita-
tional acceleration (g) acting upon a mass decreases as the mass moves away from
the Earth sea level, as follows
lE
g¼
ðRE þ hÞ2
where lE is the gravitational parameter of the Earth. The air is assumed to be dry.
The presence of humidity causes a lower density than that resulting from the tables.
However, the density decrease is usually small (less than 1%, according to Gyatt
[173]) and may be neglected.
Jacchia models have been published by Luigi G. Jacchia in 1964 [174], 1970
[175], 1971 [176], and 1977 [177]. They are also cited briefly here as J64, J70, J71,
and J77. According to Wise et al. [178], the Jacchia models use existing thermo-
spheric densities resulting from satellite drag and mass spectrometer measurements
to compute global distributions of exospheric (i.e. above about 400 km altitude)
temperature at given levels of solar flux (F10.7) and their 81-day-centred average
(F81). The choice of 81 days for the average value of solar flux depends on taking
the average on three solar Carrington rotation periods, each of which is about
27 days. These models assume that the mesopause is at 90 km altitude, where the
temperature and mass density are 183 K and 3.46 10−6 kg/m3 in J70 and J71,
and 188 K and 3.43 10−6 kg/m3 in J77. Below 90 km, atmospheric constituents
are considered to be well mixed, and the different gas species are considered to have
the same compositional ratios as at the ground. J70 and J71 only provide
mass-density profiles, whereas J77 includes profiles of two molecular and four
mono-atomic species. Above the 90-to-105 km transition layer, all gas species are
in diffusive equilibrium. Temperatures rise with altitude from the minimum at
90 km, pass through an inflection point at 125 km, and asymptotically approach a
suitable exospheric temperature (T∞) at altitudes where the mean free path for all
species becomes too long to support thermal gradients. Densities are determined for
temperature profiles associated with the assumed value of T∞, such that
dni mi g dT
¼ dz ð1 þ ai Þ
ni kB T T
where ni and mi are, respectively, the number density (number of particles per cubic
metre) and mass of the ith species; ai is the thermal diffusion coefficient of the ith
species (except for helium, for which ai = 0); g, kB, and z are, respectively, the
acceleration due to gravity, Boltzmann’s constant, and geometric altitude. The
diffusion equation written above is applied to determine local densities at any
particular altitude. It is to be noted that density results from
P the sum of the products
nimi over the number of atmospheric constituents (q = i nimi). In the model J77,
density profiles are perturbed by semi-annual variations, after the diffusion equation
written above, that is,
dni mi g dT
¼ dz ð1 þ ai Þ
ni kB T T
has been applied to determine density from a given exospheric temperature T∞.
Geomagnetic corrections enter as combinations of perturbations to both exospheric
temperatures and densities [178].
As has been shown above, the J77 model requires the value of the exospheric
temperature (T∞). This temperature varies with time and position, and must
therefore be recalculated for new evaluation of density [179].
The J71 model assumes the minimum global exospheric temperature (in the
absence of solar radiation and geomagnetic activity) to be TC = 379.0 K. This is the
night-time global exospheric temperature. The presence of solar radiation is taken
into account by means of the following expression
TC ¼ 379:0 K þ 3:24 KF81 þ 1:3 KðF10:7 F81 Þ
The last term on the right-hand side of the preceding expression takes account of
daily variations around the mean global exospheric temperature [2].
The diurnal model for T∞ also takes into account the local hour angle (H) of the
Sun with respect to the measurement point, the declination (dS) of the Sun, and the
geographic latitude (u) of the measurement point. This model is given by the

T1 ¼ TC 1 þ 0:3 sin2:2 jhj þ cos2:2 jgj sin2:2 jhj cos3:0 ðs=2Þ
where the angles s (−180° < s < 180°), h, and η are such that
s ¼ H 37 :0 þ 6 :0 sinðH þ 43 :0Þ
h ¼ 1=2ðu þ dS Þ
g ¼ 1=2ðu dS Þ
The hour angle (H) of the Sun with the measurement point is given by
H ¼ a aS
[179], a and aS being, respectively, the right ascension of the measurement point
and the right ascension of the Sun.
Montenbruck and Gill [2] note that the terms −37.0° and 6.0° sin(H + 43.0°),
which are added to H in the preceding expression
s ¼ H 37 :0 þ 6 :0 sinðH þ 43 :0Þ
account for asymmetric effects in the temperature variation relative to the position
of the Sun. They also note that the difference between the geographic latitude (u)
and the geocentric latitude (u*) is always less than 12′ and can therefore be
neglected.
In order to avoid difficulties connected with the arguments |h| and |η| of the sine
and cosine functions, Montenbruck and Gill [2] propose to use the following
trigonometric identities
sin2:2 ðjhj=2Þ ¼ ½ð1 cos hÞ=21:1

cos2:2 ðjgj=2Þ ¼ ½ð1 þ cos gÞ=21:1
in the preceding expression, which is rewritten below for convenience

T1 ¼ TC 1 þ 0:3 sin2:2 jhj þ cos2:2 jgj sin2:2 jhj cos3:0 ðs=2Þ
The correction (DT∞) to the exospheric temperature (T∞) due to the geomag-
netic activity is expressed by Jacchia [176] as follows
DT1 ¼ f DT1
H
þ ð1 f ÞDT1
L
where

H
DT1 ¼ 28:0 K Kp þ 0:03 K exp Kp
is the correction (DT∞) which applies to high altitudes (z > 350 km);

L
DT1 ¼ 14:0 K Kp þ 0:02 K exp Kp
is the correction (DT∞) which applies to low altitudes (z < 350 km);
f ¼ 1=2ftanh½0:04ðz 350 kmÞ þ 1g
is a transition function introduced by Jacchia to make the temperature correction

(DT∞) continuous at z = 350 km; Kp is the three-hour planetary geomagnetic index
for a time 6.7 h earlier than the time under consideration; and z is the geometric
altitude. Finally, the exospheric temperature (T∞), which takes account of the solar
and geomagnetic terms, is expressed by Jacchia as follows
T1 ¼ T1 þ DT1
The required value of the standard atmospheric density is computed as a function

of T∞ and z, as will be shown below.
The J71 model is based on an empirical temperature profile which starts from a
fixed value T0 = 183 K at a geometric altitude of 90 km (see above). This value
increases in the transition regime with increasing altitude and reaches asymptoti-
cally the value of T∞. Given the molecular weights and the number densities of the
atmospheric species N2, O2, Ar, and He at sea level, the standard density results
from integrating the barometric differential equation

dq T hM i hM ig
¼ d dz
q hM i T RT
[180], where q is the density, T is the temperature, ⟨M⟩ is the mean molecular mass,
R = 8.3144621 J K−1 mol−1 is the universal gas constant, z is the geometric alti-
tude, and g = 9.80665/(6356766 + z), with z in metres, is the specific force of
gravity. The barometric equation written above is used for computing densities in
the altitude range of 90–100 km.
At higher altitudes (from 100 km upwards), the diffusion differential equation
dni mi g dT
¼ dz ð1 þ ai Þ
ni kB T T
[180] is integrated, where ni, Mi, and ai are, respectively, the number density, the
molecular mass, and the thermal diffusion coefficient of the species i. Following
again de Lafontaine and Hughes [180], the mean molecular mass and density result
from
P
i ni M i
hM i ¼ P
i ni
P
ni
q ¼ hM i i
NA
where NA = 6.02214 1023 mol−1 is Avogadro’s number.

In order to integrate numerically the barometric (or diffusion) equation, a
quadrature formula can be used. This yields the number density (ni) for each
atmospheric constituent (i). This method has the advantage of requiring only a few
data to extract the full information contained in the Jacchia model J71, but also the
disadvantage of a great computational effort, because an integration is necessary
every time a density value is needed.
Another method, due to Gill [181], consists of fitting the Jacchia standard density
values, by using the method of least squares, to obtain coefficients cij to be used in
the following approximation
m X
X n
log½qðz; T1 Þ ¼ cij zi T1
j
i¼0 j¼0
In other words, given a set of k data points (zi, T∞i, qi)i=1,2,…,k in the range of
geometric altitudes going from 90 to 2500 km and in the range of temperatures
going from 500 to 1900 K, it is required to find a bi-polynomial approximation
which fits the logarithms log[q(z, T∞)] of the density values in such a way as to
have the minimum squared norm of the residuals. Let y be the following known
k 1 vector
2 3
log½qðz1 ; T11 Þ
6 log½qðz2 ; T12 Þ 7
6 7
y6 .. 7
4 . 5
log½qðzk ; T1k Þ
whose entries are the logarithms of the given k values of atmospheric density q.
Let A be the following known (m +1) (n + 1) matrix
2 3
z01 T11
0
z11 T11 1
. . . zm n
1 T11
6 z2 T12
0 0 2 1
z1 T11 . . . zm n 7
2 T12 7
6
A6 . .. .. .. 7
4 .. . . . 5
z0k T1k
0
z0k T1k
1
. . . zk T1k
m n
Let c be the following unknown (m +1) (n + 1) matrix

2 3
c00 c01 ... c0n
6 c10 c11 ... c1n 7
6 7
c 6 .. .. .. .. 7
4 . . . . 5
cm0 cm1 ... cmn
If (m +1) (n + 1) < k, it is required to determine the unknown entries cij of

c which correspond to the minimum squared norm ||q||2 = qTq of the residuals
q ¼ y Ac
To this end, as has been shown in Chap. 2, Sect. 2.10, by pre-multiplying (i.e.
by multiplying on the left) the terms on the left-hand side and on the right-hand side
of the normal equations Ac = y by the transpose (AT) of A, there results
AT Ac ¼ AT y
Then, by pre-multiplying the terms on both sides of the preceding equation by

(ATA)−1, there results
1
c ¼ AT A AT y ¼ A þ y
Gill [181] solves the problem of determining the coefficient matrix c by means of
a decomposition of the ATA matrix into the product of a lower-left triangular matrix
L and an upper-right triangular matrix U (LU-method) followed by a back-
substitution. Since it is not possible in practice to fit the full altitude and temperature
ranges by using only one approximation, Gill divides the problem into separate
sections which are solved independently. In order to avoid problems of discontinuity
in the values of density relating to adjacent sections, Gill uses a constrained
least-squares polynomial fit, which ensures continuity in densities and in their first
derivatives with respect to altitude at the boundary of each section. By so doing, he
finds that a fifth-order polynomial in geometric altitude (z) and a fourth-degree
polynomial in exospheric temperature (T∞) lead to satisfactory results, in case of
boundaries placed at, respectively, 850 K, 180, 500, and 1000 km.
In other words, Gill sets m = 5 and n = 4. The results are expressed in form of
two sets of tables (courtesy of Gill [181]), each of which sets comprises four tables.
The first set of tables, given below, refers to the standard density polynomials in
temperature (index j) and altitude (index i) below 500 km.
90 km < z < 180 km, 500 K < T∞ < 850 K
i\j 0 1 2 3 4
0 −0.352085629 102 +0.391262208 101 −0.864925858 102 +0.150411921 103 −0.710942779 102
1 +0.112921028 104 +0.119815810 104 +0.863379360 103 −0.357709109 104 +0.197055758 104
2 −0.152747516 10 5
−0.355848094 10 5
+0.189924264 10 5
+0.250824047 105
−0.196825320 105
3 +0.930204177 105 +0.364655435 106 −0.329036395 106 −0.120963117 105 +0.843813722 105
4 −0.273439442 106 −0.157609709 107 +0.168583089 107 −0.428294274 106 −0.134559332 106
5 +0.314969634 10 6
+0.248772280 107
−0.289912393 10 7
+0.111190398 107
+0.329409498 104
90 km < z < 180 km, 850 K < T∞ < 1900 K
i\j 0 1 2 3 4
0 −0.533541178 102 +0.290055675 102 −0.204643854 102 +0.797714877 101 −0.133585294 101
1 +0.197753256 104 −0.709147804 103 +0.439853751 103 −0.156871952 103 +0.261546624 102
2 −0.299362006 105 +0.518728595 104 −0.198979486 104 +0.364316576 103 −0.570066884 102
3 +0.211206780 106 −0.448302904 104 −0.134997048 105 +0.951001182 104 −0.165372508 104
4 −0.720972149 10 6
−0.768410087 10 5
+0.125623620 10 6
−0.680569856 10 5
+0.118125702 105
5 +0.962596639 106 +0.212312699 106 −0.262279285 106 +0.133712984 106 −0.232999509 105
180 km < z < 500 km, 500 K < T∞ < 850 K
i\j 0 1 2 3 4
0 +0.231191014 102 +0.135529793 103 −0.842431012 103 +0.128733075 104 −0.618120934 103
1 −0.105777591 104 +0.608797320 103 +0.869056563 104 −0.171592213 105 +0.905267060 104
2 +0.117722961 105 −0.316413234 105 −0.107632268 104 +0.630262927 105 −0.431245935 105
3 −0.582766334 10 5
+0.218816743 106
−0.242291149 10 6
+0.246128588 105
+0.604409552 105
4 +0.125458876 106 −0.543470990 106 +0.812301589 106 −0.449043841 106 +0.500745829 105
5 −0.945292227 105 +0.440802570 106 −0.737940971 106 +0.509527260 106 −0.115419197 106
180 km < z < 500 km, 850 K < T∞ < 1900 K
i\j 0 1 2 3 4
0 +0.404176122 102 −0.130571901 103 +0.146680893 103 −0.712029554 102 +0.126960539 102
1 −0.812772013 103 +0.227356511 104 −0.257726141 104 +0.125904483 104 −0.225497799 103
2 +0.513004248 10 4
−0.150130763 10 5
+0.171714170 10 5
−0.844169798 10 4
+0.151879634 104
(continued)
(continued)
i\j 0 1 2 3 4
3 −0.160017007 105 +0.477046885 105 −0.547349228 105 +0.269966778 105 −0.487030568 104
4 +0.238471768 105 −0.719906407 105 +0.828465304 105 −0.409835788 105 +0.741192600 104
5 −0.136310388 105 +0.415349930 105 −0.479358081 105 +0.237785404 105 −0.431023262 104
Likewise, the second set of tables, given below, refers to the standard density
polynomials in temperature (index j) and altitude (index i) above 500 km.
500 km < z < 1000 km, 500 K < T∞ < 850 K
i\j 0 1 2 3 4
0 −0.181572147 104 +0.979297163 104 −0.183137427 105 +0.138525479 105 −0.345123396 104
1 +0.985122122 104 −0.539752453 105 +0.999316898 105 −0.725945601 105 +0.162255300 105
2 −0.182293162 105 +0.100242982 106 −0.178448116 106 +0.114517833 106 −0.164193391 105
3 +0.129811344 10 5
−0.711343013 10 5
+0.110637478 10 6
−0.382577652 10 5
−0.166691482 105
4 −0.153350954 104 +0.781553694 104 +0.703756219 104 −0.467463592 105 +0.351694860 105
5 −0.126367969 104 +0.726579206 104 −0.209290906 105 +0.293609410 105 −0.149167582 105
500 km < z < 1000 km, 850 K < T∞ < 1900 K
i\j 0 1 2 3 4
0 −0.402133495 102 −0.132698258 103 +0.377886396 103 −0.280866017 103 +0.651353118 102
1 +0.425578905 103 +0.352812623 103 −0.207788843 104 +0.172654309 104 −0.419147661 103
2 −0.182166204 10 4
+0.790535705 103
+0.393427057 10 4
−0.396933359 10 4
+0.102799073 104
3 +0.307023080 104 −0.294154027 104 −0.327663903 104 +0.442021654 104 −0.123077830 104
4 −0.219684753 104 +0.258511809 104 +0.138277623 104 −0.253300602 104 +0.745138710 103
5 +0.549495863 10 3
−0.660422507 10 3
−0.332807663 10 3
+0.633570336 103
−0.187981171 103
1000 km < z < 2500 km, 500 K < T∞ < 850 K
i\j 0 1 2 3 4
0 +0.354869756 103 −0.250868544 104 +0.625274219 104 −0.675537591 104 +0.267576301 104
1 −0.537085151 103 +0.418258598 104 −0.115111424 105 +0.133891525 105 −0.561058041 104
2 −0.234958642 102 −0.894184106 103 +0.441792700 104 −0.673281743 104 +0.331260819 104
3 +0.340707342 103 −0.153158821 104 +0.217904521 104 −0.884134055 103 −0.136976861 103
4 −0.169847098 10 3
+0.898569661 103
−0.170479728 10 4
+0.136309843 104
−0.381241735 103
5 +0.249797337 102 −0.138961772 103 +0.282005813 103 −0.247286186 103 +0.789643901 102
1000 km < z < 2500 km, 850 K < T∞ < 1900 K
i\j 0 1 2 3 4
0 +0.128106144 102 −0.338917904 103 +0.686193462 103 −0.466762722 103 +0.102966235 103
1 +0.202425105 103 +0.166830201 103 −0.114787559 104 +0.991894020 103 −0.243021528 103
2 −0.575074276 103 +0.825982311 103 +0.232983196 103 −0.650335882 103 +0.199798890 103
3 +0.510620714 103 −0.103201223 104 +0.485187352 103 +0.821409740 102 −0.652704815 102
4 −0.189895280 103 +0.434750107 103 −0.298601120 103 +0.542317977 102 +0.503945902 101
5 +0.256957682 10 2
−0.628271027 10 2
+0.497107699 10 2
−0.140438474 10 2
+0.845050002 100
The approximation found by Gill holds in the altitude range going from 90 to
2500 km and in the temperature range going from 500 to 1900 K. After computing
the standard value of density, or rather its logarithm log[q(z, T∞)], by using the Gill
polynomials, it is necessary to apply several corrections, in order to take account of
the effects which have been mentioned above. First, at an altitude less than 350 km,
it is necessary to add the geomagnetic (subscript GM) term
D log qGM ¼ ½0:012 Kp þ 1:2 105 expðKp Þð1 f Þ
to the logarithm of the standard value of density, where Kp is the geomagnetic

index, and f is the transition function of the geometric altitude z expressed by
f ¼ 1=2ftanh½0:04ðz350 kmÞ þ 1g
Second, the semi-annual (subscript SA) density variation in the thermosphere

and in the lower exosphere is given by the following empirical expression
D log qSA ¼ f ðzÞgðtÞ
where g(t) is the temporal variation and f(z) is the amount of density variation at a
given geometric altitude. Jacchia obtained these functions as best fits of the actual
data made available by satellites, as follows
f ðzÞ ¼ ½5:876 107 z2:331 þ 0:06328 expð0:002868zÞ

gðtÞ ¼ 0:02835 þ ½0:3817 þ 0:17829 sinð2psSA þ 4:137Þ sinð4psSA þ 4:259Þ
where z is the geometric altitude expressed in km, t is the time expressed in

Modified Julian Days (MJD = JD−2400000.5), and
( 1:65 )
1 1 1
sSA ¼ U þ 0:09544 þ sinð2pU þ 6:035Þ
2 2 2
t 36204
U¼
365:2422
U is the number of tropical years elapsed since the 1st of January 1958. By the way,
a tropical year is the time which the Sun takes to return to the same position in the
cycle of seasons as seen from the Earth (e.g. the time between two consecutive
vernal equinoxes or two consecutive summer solstices). The maximum value of the
logarithmic semi-annual density correction is about 0.21.
Third, the seasonal-latitudinal (subscript SL) density variation results from
sin3 u
D log qSL ¼ 0:014Dz90 exp 0:0013Dz290 sinð2pU þ 1:72Þ
jsin uj
where Dz90 = z − 90 km, and u is the geographic latitude.
In this regard, Montenbruck and Gill [2] observe that, in computer applications,
the expression sin3u/|sin u| should be replaced by SIGN(sin2u, u).
The fourth cause of density variation is a strong increase in the concentration of
Helium above the Earth winter pole, that is, above the north pole in December and
above the south pole in July. The Jacchia 1971 model takes account of this phe-
nomenon by means of an empirical expression, as will be shown below.
A polynomial approximation, like that which is used to compute the standard
values of atmospheric density, has been proposed by Gill in order to compute the
atmospheric density variations due to Helium.
The correction DqHe to the standard atmospheric density due to variations in
Helium density results from the following expressions

dS 3 p u dS

Dv D logðnHe Þ ¼ 0:65 sin 0:35355
4 2 j dS j
m X
X n
v log½nHe ðz; T1 Þ ¼ hij zi T1
j
i¼0 j¼0
mHe Dv
DqHe ¼ 10v 10 1
NA
where NA = 6.02214 1023 mol−1 is Avogadro’s number, mHe = 4.0026 is the

molecular weight of Helium, dS is the declination of the Sun, e is the obliquity of
the ecliptic, z is the geometric altitude, T∞ is the exospheric temperature, and hij are
the coefficients which appear in the series expansion of m. These coefficients have
their index of height i ranging from 0 to m = 5 and their index of temperature
j ranging from 0 to n = 4, as is the case with the coefficients cij considered above for
the standard values of atmospheric density. The values of the coefficients hij are
given in the following tables (courtesy of Gill [181]).
90 km < z < 500 km
i\j 0 1 2 3 4
0 +0.183154936 10+2 +0.588755629 10−2 −0.481325706 10−5 +0.170173800 10−8 −0.212837367 10−12
1 −0.737400840 10−1 −0.125107697 10−3 +0.103926920 10−6 −0.367927971 10−10 +0.455525822 10−14
2 +0.438416419 10−3 +0.865702730 10−6 −0.721694577 10−9 +0.248153370 10−12 −0.285907444 10−16
3 −0.141119464 10−5 −0.248383435 10−8 +0.200410722 10−11 −0.624498540 10−15 +0.556100355 10−19
4 +0.215363863 10−8 +0.342194382 10−11 −0.262896070 10−14 +0.708565525 10−18 −0.327980417 10−22
5 −0.125513854 10−11 −0.182725298 10−14 +0.132158108 10−17 −0.288739764 10−21 −0.782717842 10−26
500 km < z < 1000 km
i\j 0 1 2 3 4
0 +0.162708865 10+2 −0.178681641 10−1 +0.307907898 10−4 −0.204343061 10−7 +0.464341888 10−11
1 −0.195829652 10−1 +0.138612645 10−3 −0.253246272 10−6 +0.171418298 10−9 −0.393423008 10−13
2 +0.251425106 10−4 −0.380633890 10−6 +0.769237606 10−9 −0.539476570 10−12 +0.126030389 10−15
−7 −9 −11 −15
3 −0.298331353 10 +0.585585083 10 −0.121066317 10 +0.856163162 10 −0.200903018 10−18
4 +0.180202809 10−10 −0.438287770 10−12 +0.920153031 10−15 −0.654393514 10−18 +0.154021961 10−21
−14 −15 −18 −21
5 −0.424306744 10 +0.126883016 10 −0.269580744 10 +0.192546917 10 −0.454232921 10−25
1000 km < z < 2500 km
i\j 0 1 2 3 4
0 +0.187334592 10+2 +0.228568321 10−1 −0.686077568 10−4 +0.537962260 10−7 −0.132755885 10−10
1 −0.236253018 10−1 −0.690761294 10−4 +0.225167974 10−6 −0.179593742 10−9 +0.446365926 10−13
2 +0.189389910 10−4 +0.114595980 10−6 −0.318325867 10−9 +0.246107555 10−12 −0.604042307 10−16
3 −0.113219809 10−7 −0.743832646 10−10 +0.204028767 10−12 −0.157319125 10−15 +0.385703179 10−19
−11 −13 −16 −19
4 +0.346501347 10 +0.230894316 10 −0.632046641 10 +0.487141850 10 −0.119413886 10−22
5 −0.415670958 10−15 −0.279192982 10−17 +0.763279237 10−20 −0.588111152 10−23 +0.144145505 10−26
Another atmospheric model used frequently in the scientific community is the

NRLMSISE-00, where NRL stands for Naval Research Laboratory of the United
States Navy, MSIS stands for Mass Spectrometer and Incoherent Scatter Radar, and
E stands for extended, because this model extends from the ground to space, and 00
stands for the year of release of the model. It is fully described in Ref. [184]. For the
purposes of this book, the main properties of this model are given in the following
figure (courtesy of Wikimedia [186]), which shows graphically the air density
computed by the NRLMSISE-00 as a function of altitude.
The NRLMSISE-00 is an empirical atmospheric model, whose field of extension

goes from the surface of the Earth to the upper limit of the thermosphere (about
600 km from the ground). It is the upgrade of a previous model named MSIS-E-90.
These models use data on the atmospheric composition obtained by satellites and
temperatures coming from ground-based radars.
In particular, the NRLMSISE-00 model is based on:
• data on atmospheric drag measured in operations of orbit determination;
• further data of the same type measured by satellite-borne accelerometers;
• data on temperatures resulting from incoherent scatter radar observations per-
formed at the Millstone Hill and Arecibo observatories; and
• observations of molecular oxygen (O2) performed by the Solar Maximum
Mission and based on solar ultraviolet occultation.
In addition, recent investigations of atmospheric drag measured by satellites in
low-Earth orbit have shown that the Jacchia model J70 attributes erroneously
non-thermospheric drag sources to atomic helium (He). By contrast, Picone et al.
[182] point out that such sources are “hot” atomic oxygen and ionospheric atomic
oxygen ions (O+). Since none of these species is in thermal equilibrium with the
thermosphere, the NRLMSISE-00 model treats them as a new component, called
“anomalous oxygen”, to drag.
This model has a web interface [184], which requires the user to specify the
following values: date (year, month and day), time (universal or local), geographic
or geomagnetic co-ordinates (latitude and longitude), altitude (from 0 to 1000 km),
and optionally the index of geomagnetic activity (ap), and the solar proxy (F10.7).
On the basis of such data, the model computes the following values: number
densities of atmospheric constituents (including anomalous oxygen), total mass
density, exospheric temperature, and temperature at altitude.
Reference [185] tabulates the concentrations, temperatures, and densities as a
function of altitude for the NRLMSISE-00 model. The data are given in the range
of altitudes going from ground level up to 900 km, at equatorial latitude, and are
contained in three tables: Table 3.1 (pp. 20–21) refers to low solar and geomagnetic
activities, Table 3.2 (pp. 22–23) refers to mean solar and geomagnetic activities,
and Table 3.3 (pp. 23–24) refers to high long-term solar and geomagnetic activities.
According to Doornbos [165], the main advantage which comes from using
datasets based only on mass spectrometer and incoherent scatter radar observations
instead of drag-derived datasets is that the former datasets consist of independent
observations of both temperature and number densities for the atmospheric
constituents.
Still another atmospheric model is the Jacchia-Bowman 2008 (also known as
JB2008) model, which is an improvement of the Jacchia models J64, J70, J71, and
J77 described above, and gives the neutral temperature and the total density in the
thermosphere and exosphere above an altitude of 120 km. This model is fully
described in Ref. [185]. Its main properties are given below. The authors of the
JB2008 model use further solar indices (or proxies) than the 10.7-cm solar flux
(F10.7) to compute the value of the global night-time minimum exospheric tem-
perature (TC). By the way, Tobiska et al. (186) distinguish between solar irradiance
indices and proxies. According to their definition, a solar irradiance proxy is a
measured or modelled data type that is used as a substitute for solar spectral
irradiances, whereas a solar irradiance index is a measured or modelled data type
that is an indicator of a solar spectral irradiance activity level.
As has been shown above, the J71 model assumes the minimum global exo-
spheric temperature (in the absence of solar radiation and geomagnetic activity) to
be TC = 379.0 K. This is the night-time global exospheric temperature. In the J71
model, the presence of solar radiation is taken into account by means of the fol-
lowing expression
TC ¼ 379:0 K þ 3:24 K F81 þ 1:3 K ðF10:7 F81 Þ
where F81 is the 81-day-centred average of the solar flux F10.7, and the last term on
the right-hand side of the preceding expression takes account of daily variations
around the mean global exospheric temperature [2].
By contrast, the JB2008 model computes the value of TC by means of the
following expressions

14
F81
WT ¼
240
FS ¼ F81 WT þ S81 ð1 WT Þ
TC ¼ 392:4K þ 3:227KFS þ 0:298KðF10:7 F81 Þ þ 2:259KðS10:7 S81 Þ
þ 0:312KðM10:7 M81 Þ þ 0:178KðY10:7 Y81 Þ
The solar proxy F10.7 and its 81-day-centred average F81 have been described
above. We only add here that the observed F10.7 values and their 81-day-centred
average F81, with a one-day lag, are used in the JG2008 model.
The new indices (or proxies) appearing in the preceding expressions are S10.7,
M10.7, and Y10.7, and their respective 81-day-centred averages S81, M81, and Y81.
The meaning of these new quantities is explained below. Following Tobiska (187),
the S10.7 index gives the integrated 26–34 nm solar irradiance measured by the solar
extreme-ultraviolet monitor (SEM) instrument on the NASA/ESA Solar and
Heliospheric Observatory (SOHO) satellite. This satellite moves around the Sun at
the same speed as the Earth, by slowly orbiting around the First Lagrangian Point
(L1), where the combined gravitational forces exerted by the Earth and by the Sun
on the satellite keep it in an orbit locked to the Earth–Sun line. The L1 point is about
1.5 106 km away from the Earth, in the direction of the Sun. This particular orbit
has been chosen in order for the satellite to have an uninterrupted view of the Sun.
The SEM instrument has made observations since the 16th of December 1995 for
the 26–34 nm solar extreme ultraviolet (EUV) emission with a time resolution of
15 s. The data gathered by the instrument are first normalised, by dividing the daily
value by a mean value of 1.9955 1010 photons cm−2 s−1 for the time interval
going from 16 December 1995 to 12 June 2005 (solar cycle 23), and then converted
to SFU units by using a first-degree polynomial fit with F10.7, in such a way as to
exclude spikes from abnormal flares and missing data. This done, the value of the
S10.7 index to be used in the JB2008 model results from the following expression

SOHO SEM2634
S10:7 ¼ 2:90193 þ 118:512
1:9955 1010
The M10.7 proxy comes from the Mg II core-to-wing ratio measurements per-
formed by operational satellites of the NOAA series, which carry the solar
backscatter ultraviolet (SBUV) spectrometer, aimed at monitoring ozone in the
lower atmosphere of the Earth, and also by the research satellites SORCE/
SOLSTICE and ERS-2/GOME. The Mg II core-to-wing ratio is computed by
taking the ratio of the intensity of the Mg II emission at 280 nm to that of the
nearby line wings. This is done in order to make the measurements of the solar
activity in the chromospheric region and in some photospheric regions independent,
as much as possible, of changes of instrument sensitivity with time. This ratio is a
good proxy for some solar energy emissions in the far and in the extreme regions of
the ultraviolet spectrum. The Mg II core-to-wing ratio is used in a linear regression
with F10.7 to derive the M10.7 proxy in SFU units, as follows. Let MgIIcwr_SET be the
data product (Mg II core-to-wing ratio) available at at the Products menu link of the
website http://spacewx.com. In order to determine the value of M10.7 to be used in
the JB2008 model, it is necessary to compute first M*10.7, which is the Mg II
core-to-wing ratio scaled to F10.7 in SFU units, and then

M10:7 ¼ 2107:6186 þ ð8203:0537Þ MgIIcwr SET þ f M10:7 ½1:2890589
þ ð8:3777235 105 Þ x 1g
where x = 0, 1, 2, … is the number of days elapsed from 2448542.0 JD (that is,

from 12 October 1991, 12:00 UT). A two-day lag time is appropriate for M10.7.
The Y10.7 index is a composite quantity meant to represent the 0.1–0.8 nm X-ray
emission during solar maximum activity and the Lyman-a emissions during mod-
erate and low solar activities. By the way, the Lyman-a line (whose wavelength k is
121.6 nm) is a spectral line of hydrogen emitted when the electron falls from the
n = 2 orbital to the n = 1 orbital, where n is the principal quantum number.
Therefore, the Y10.7 index takes account of the combined energy input resulting
from the solar Lyman-a radiation, which is emitted by the chromospheric/transition
region, and also from the hot coronal X-rays. The radiation of both of these types
comes from the Sun to the Earth and penetrates to the lower thermosphere and
mesopause (85–100 km). There, the radiation is absorbed by molecular oxygen
(O2) and molecular nitrogen (N2) and ionises such neutral constituents. The data on
the 0.1–0.8 nm X-rays are measured by the X-ray Spectrometer (XRS) instrument
carried by the GOES series operational spacecraft, where GOES stands for
Geostationary Operational Environmental Satellites. Such data are continuously
reported by the Space Weather Prediction Centre (SWPC) of the National Oceanic
and Atmospheric Administration (NOAA) in the website http://www.swpc.noaa.
gov/. Tobitska and Bouwer have created an index, called Xb10, based on these data
for the purpose of representing the daily energy absorbed by the mesosphere and
lower thermosphere. The data on the Lyman-a radiation are regularly measured by
the SOLSTICE (SOLar STellar Irradiance Comparison Experiment) instrument on
the UARS (Upper Atmosphere Research Satellite) and SORCE (Solar Radiation
and Climate Experiment) satellites. Such data are also measured by the SEE (Solar
EUV Experiment) carried by the TIMED (Thermosphere Ionosphere Mesosphere
Energetics and Dynamics) spacecraft. The composite Y10.7 index, in SFU units, is
based on the Xb10 index, which represents the X-rays, and on a further index, called
Lya, which represents the Lyman-a radiation, as shown by the following
expressions
X10 ¼ 42:5991 þ 0:533669 Xb10

L10 ¼ 88:3926 þ 3:35891 1010 Lya þ 2:40481 1022 Lya2
Y10:7 ¼ F81norm X10 þ ð1 F81norm Þ L10
where F81norm is the normalised F81, that is, the 81-day-centred smoothed solar
index F10.7 divided by its mean value in the common time interval going from the
1st of January 1991 to the 16 th of February 2008, and X10 has a minimum threshold
value of 40. The Y10.7 index is used in the JB2008 model with a 5-day lag.
Reference values of the four solar indices or proxies described above, for
intermediate-term and short-term solar variability, have been given by Tobitska
[187]. In addition, daily and 81-day averaged values of such indices or proxies are
given by Space Environment Technologies (SET) in the file SOLFSMY.TXT,
updated daily, which can be found at the website http://sol.spacenvironment.net/
*JB2008/indices.html.
The geomagnetic activity is represented in the JB2008 model by two indices, for
the purpose of modelling the low, unsettled, and the sub-storm/storm activity. Such
indices are ap and Dst. The former index (ap) has been described above. When the
value of ap is less than or equal to 40, then the JB2008 model does not use (as is the
case with the Jacchia models) the value of Dst, because the geomagnetic activity is
assumed to be low or unsettled. In such cases, the correction to the exospheric
temperature due to the geomagnetic activity is computed in the JB2008 model in
the same way as that shown above for the Jacchia models, that is, by using the value
of Kp, which is the three-hour planetary geomagnetic index for a time 6.7 h earlier
than the time under consideration. By contrast, for values of ap greater than 40, the
model assumes a geomagnetic storm or sub-storm to be in progress and, therefore,
does use the value of Dst. By the way, a geomagnetic storm is a temporary dis-
turbance of the magnetic field of the Earth.
Dst stands for Disturbance Storm Time. The value of this index, expressed in
nanoteslas (nT), with 1 nT = 1 10−9 T = 1 10−9 Wb/m2, measures the
strength of geomagnetic storms. This value, if negative, indicates the presence of a
geomagnetic storm at the time of measurement, with higher negative values of Dst
corresponding to stronger geomagnetic storms.
These increasing negative values are caused by the ring current flowing tor-
oidally westward (i.e. clockwise if viewed by an observer placed in the northern
hemisphere) at altitudes of about 10000–60000 km in the equatorial plane of the
Earth, as shown in the preceding figure, which is due to the courtesy of the United
States Air Force [188].
According to Daglis et al. [189], the main carriers of the terrestrial ring current
are positive ions, with energies ranging from about 1 keV to a few hundred keV,
which are trapped by the geomagnetic field and undergo an azimuthal drift. Of
course, positive ions and negative electrons drift in opposite directions. This flow of
charged particles produces a current, whose magnetic field opposes to the geo-
magnetic field. Therefore, the geomagnetic field decreases in this area.
The Dst index is based on the average value of the horizontal component of the
geomagnetic field measured hourly at four near-equatorial geomagnetic observa-
tories. Real-time values of the Dst index are made available by the World Data
Centre (WDC) for Geomagnetism (Kyoto, Japan) through the website http://wdc.
kugi.kyoto-u.ac.jp/dst_realtime/index.html. Values of the Dst index, updated daily,
are also given in the file DSTFILE.TXT available at the website http://sol.
spacenvironment.net/*JB2008/indices.html.
According to Tsurutani and Gonzalez [190], a geomagnetic storm is composed

of three phases. In the first (or initial) phase, the horizontal component of the
geomagnetic field increases to positive values of up to tens of nanoteslas. In the
second (or main) phase, the same component decreases to minus hundreds of
nanoteslas. In the third (or recovery) phase, the geomagnetic field returns gradually
to its quiet time value. The duration of these three phases is variable. The initial
phase may last from some minutes to many hours; the main phase, from half an
hour to several hours; and the recovery phase, from tens of hours to a week. In the
presence of a geomagnetic storm or sub-storm (ap > 40), the JB2008 model
computes the correction to the exospheric temperature due to the geomagnetic
activity by using the Dst index, as described in detail by Bowman and Wolfe [191].
A brief account of the method used by them to compute the correction is given
below. Burke et al. [192] have shown that, in case of a geomagnetic storm, the
thermosphere acts as a driven-but-dissipative thermodynamic system, which is
mathematically described by the following first-order differential equation

1 aT 1
T1SW ðtn þ 1 Þ ¼ 1 T1SW ðtn Þ þ Dstðtn þ 1 Þ 1 Dstðtn Þ
sE aD sRC
where T∞ is the orbit-averaged exospheric temperature, the subscript SW denotes

the contribution to T∞ due to the solar wind, tn and tn+1 are two successive times, sE
is the relaxation time constant of the global energy, whose value (6.5 h) is deter-
mined empirically, aT and aD are coupling constants, whose ratio aT/aD has been set
equal to about 1.575 by Burke et al. in order to integrate numerically the preceding
equation, sRC is the relaxation time constant of the ring current (subscript RC),
whose value has been set equal to 7.7 h by Burke et al., and Dst(tn) and Dst(tn+1)
are the values of the Dst index at the times, respectively, tn and tn+1. The differential
equation given above must be integrated numerically from the initial time to the
final time of any geomagnetic storm, in order to compute T∞SW at every time during
the main and early recovery phases of storms. In the numerical integration, Burke
et al. have assumed T∞SW(t0) = 0. In addition, they have assumed that the Dst index
has its actual value if this value is less than −20 nT, whereas the value of Dst is set
equal to −20 nT in any other case, that is,
Dst ¼ Dst ðif Dst\ 20 nTÞ

Dst ¼ 20 nT ðif Dst
20 nTÞ
The values of the time constants (sE and sRC) used in the preceding differential
equation have been determined by Burke et al. on the basis of the data measured by
the accelerometers carried by the CHAMP (CHAllenging Minisatellite Payload)
and GRACE (Gravity Recovery And Climate Experiment) satellites.
Other values (e.g. sE = ∞ h and sRC = 1.0 h) for the same time constants have
also been used for the late recovery phases of storms.
The preceding figure (courtesy of COSPAR [162]) shows the mean air density
computed by the JB2008 model as a function of altitude for low, moderate, and
high long-term and short-term solar and geomagnetic activity.
3.23 The Angular Velocity of the Atmosphere
In the preceding sections, the atmosphere has been assumed to rotate with the
angular velocity of the solid surface of the Earth, due to the friction between this
surface and the atmosphere. In other words, the ratio
xatm
K¼
xE
of the angular velocity of rotation of the atmosphere to that of the solid Earth has
been taken equal to unity and constant at any altitude. However, some authors
[193–195] have shown that this assumption is not always justified.
In particular, King-Hele (193) has shown that the upper atmosphere, at heights
ranging from 200 to 350 km, rotates on average faster than the Earth, and that the
average rate of rotation increases with height from about 1.1 rev/day at 200 km to
nearly 1.4 rev/day at 350 km. This rotation rate decreases above 350 km, to about
1.0 rev/day at 420 km and 0.7 rev/day at 500 km. In addition, at heights ranging
from 120 to 230 km, there are wide variations in the rotation rate over short
intervals of time. These results are based on the reduction of the inclination of the
orbital plane of some artificial satellites with respect to the equatorial plane of the
Earth, due to the bi-normal component (m aDh) of the aerodynamic force (m aD), as
has been shown in detail in Sect. 3.20. This force reduces continuously but very
slowly the inclination of the orbital plane.
In order to take account of the actual value of K, the expression given in
Sect. 3.19 for the perturbing acceleration, aD, acting on a satellite due to atmo-
spheric drag can be written as follows
1 CD A
aD ¼ q jvrel jvrel
2 m
where vrel = v − vatm is the velocity vector of the satellite with respect to the
atmosphere, v is the velocity vector of the satellite in the GCRS (see Sect. 3.8), vatm
is the velocity vector of the atmosphere in the GCRS, q is the atmospheric density
at the point of interest, and CD A/m is the ballistic coefficient of the satellite.
The velocity vector of the atmosphere can be written, in terms of geocentric
inertial quantities, as follows
vatm ¼ xatm r
where xatm is the angular velocity vector of the atmosphere in the GCRS and r is
the position vector of the satellite in the GCRS (see Sects. 3.10 and 3.11).
Let xE be the true angular rotation vector of the Earth in the GCRS, determined
as has been shown in Sect. 3.6. In order to take account of the actual value of K
(which is not, in the general case, equal to unity), the angular velocity vector of the
atmosphere in the GCRS can be determined as follows
xatm ¼ KxE
When the six orbital elements are known at an initial epoch, t0, and the equations
of motion to be integrated are the Lagrangian equations in Gaussian form, that is,
2a2 h p i
h r
1
h

r cosðx þ /Þ
i0 ¼ ah
h

r sinðx þ /Þ
X0 ¼ a/
h sin i
3.23 The Angular Velocity of the Atmosphere 543

0 p cos / ðp þ r Þ sin / r sinðx þ /Þ cos i
x ¼ ar þ a/ ah
he he h sin i
b
M0 ¼ n þ ½p cos / 2er ar ½ðp þ r Þ sin /a/
ahe
then the three components (radial, transverse, and bi-normal) of the perturbing
acceleration vector, aD, due to drag are respectively

12
1 CD A l
aDr ¼ q jvrel j e sin /
2 m p
"
1 #
1 CD A l 2
aD/ ¼ q jvrel j ð1 þ e cos /Þ KxE r cos i
2 m p
1 CD A
aDh ¼ q jvrel jKxE r cosðx þ /Þ sin i
2 m
where p = a(1 − e2) is the semi-latus rectum, r = p/(1 + e cos /) is the radius
vector, h = (l/p)½ is the moment of momentum per unit mass, n = (l/a3)½ is the
mean motion,
8 "
1 #2 91
<l l 2 =2
2
jvrel j ¼ e sin / þ
2 2
ð1 þ e cos /Þ KxE r cos i þ ½KxE r cos u sin i
:p p ;
is the magnitude of the velocity vector of the satellite with respect to the atmo-
sphere, and u = x + / is the argument of latitude at epoch.
3.24 The Relativistic Perturbations
The purpose of the present section is to show the relativistic perturbations in the
motion of an artificial satellite orbiting around the Earth.
Le Verrier, the French astronomer whose fame rests on the prediction of the
existence of Neptune on the basis of perturbations observed on the orbit of Uranus,
was the first to report in 1859 [196] that the residual slow (about 43 arcseconds in a
century, according to Bogorodskii [197]) precession of the orbit of Mercury around
the Sun could not be explained by perturbations due to other planets known at that
time.
The precession of the perihelion of Mercury is shown in the preceding figure,

due to the courtesy of NASA (198), where the amount of perihelion shift is highly
exaggerated for the sake of clarity.
Following Odenwald (199), the total precession of Mercury’s perihelion, mea-
sured from the Earth, is 5599.74 ± 0.41 arcseconds/century. This total amount
includes the following contributions:
5025.64 ± 0.50″/century due to the precession of the equinoxes (see Sect. 3.3)
531.54 ± 0.68″/century due to gravitational forces of other planets
42.56 ± 0.94″/century due to general relativity and solar oblateness
As will be shown below, the only term having to do with the general relativity
theory is the last term, which is mainly due to the curvature of space-time in the
vicinity of Mercury’s orbit.
Le Verrier considered the timed observations available at that time (from 1697 to
1848) of the transits of Mercury across the disc of the Sun. By so doing, he computed
a perihelion advance of 527 arcseconds per century due to the perturbing gravita-
tional forces exerted by the other planets known at that time. He suggested, among
possible explanations of the discrepancy observed by him (38 arcseconds/century)
between the measured and the predicted shift, that either an unknown planet or rather
a series of corpuscles might be placed in an orbit closer to the Sun than that of
Mercury. In 1895, Newcomb repeated Le Verrier’s calculations and found a shift of
43 arcseconds/century in Mercury’s perihelion, which value was higher than that
found by Le Verrier, and could not be explained. Neither the supposed planet nor the
asteroid belt between Mercury and the Sun was ever found, and the advance of
Mercury’s perihelion was eventually explained in 1915 (200) by Einstein’s general
relativity theory, as will be shown below. According to Einstein, space and time are
combined one with the other into a single four-dimensional manifold (or continuum),
called space-time, such that three of its four dimensions are the three space
co-ordinates, and the fourth dimension is time. The mass of an object creates a
concavity in space-time, in correspondence of that object, as shown in the following
figure, due to the courtesy of NASA (201). The gravitational force is the curvature of
space-time. In other words, the mass or the energy possessed by an object causes a
curvature of space-time. A moving mass or a propagating radiation must follow the
3.24 The Relativistic Perturbations 545
shortest path in space-time. In a flat two-dimensional space (that is, in a plane), the
shortest path between two points is a straight line. On the contrary, on the surface of a
sphere, the shortest path is not a straight line, but a great circle.
Generally speaking, in a positively curved space (as is the case with the surface
of a sphere) or in a negatively curved space (as is the case with the surface of a
saddle), shortest paths are not straight lines.
In space-time, geodetics are curved around massive objects. A geodetic is the

path of shortest distance between two points in space. Beams of light follow
geodetics in space-time. Since geodetics are curved by the mass of the Sun, so does
a beam of light travelling along a geodetic, and therefore the geometry of
space-time must also be curved around the Sun. Following Bogorodskii [197], the
residual secular advance of Mercury’s perihelion, due to relativistic effects, is
quantitatively expressed by the following formula
6plS
Dx ¼
c 2 að 1
e2 Þ
where Dx is the secular perturbation, in radians, of the argument of perihelion per

revolution of Mercury, lS is the gravitational parameter of the Sun, c is the speed of
light, and a and e are, respectively, the major semi-axis and the eccentricity of the
orbit of Mercury. By introducing the following values
lS ¼ 132712440000 km3 =s2

c ¼ 299792:458 km=s
a ¼ 57909100 km
e ¼ 0:20563
in the formula given above, we have
6 132712440000
Dx ¼ p
299792:4582 57909100 ð1 0:205632 Þ
¼ p 1:59748865566 107 radians/revolution
Since a revolution of Mercury takes
87:969
87:969 days ¼ centuries
365:242 100
then the same result can be expressed in arcseconds/century as follows
180 36524:2
Dx ¼ p 1:59748865566 107 3600 ¼ 42:9800 =century
p 87:969
By the way, the value computed above for Dx is the same as that computed by
Pireaux and Rozelot (202) and does not take into account the flattening of the Sun,
that is, the solar quadrupole moment (J2). By also taking into account the quad-
rupole moment and using a value of 2 10−7 for J2, Pireaux and Rozelot obtain a
value in the range going from 43.00″/century to 43.01″/century for Dx.
In 1918, about three years after the fundamental work of Einstein, Lense and
Thirring (203) discovered that a slowly rotating spherical mass generates a gravi-
tational field which is not the same as that which would be generated by a stationary
spherical mass but rather drags the space-time around it. It follows that the spinning
Earth not only curves (geodetic effect) but also twists (frame-dragging effect) the
space-time, pulling it around into a four-dimensional swirl, as shown in the fol-
lowing figure (courtesy of the Stanford University [204]).
The frame-dragging effect generates further perturbations than those due to the
geodetic effect on the motion of bodies orbiting around a rotating mass. The
3.24 The Relativistic Perturbations 547
frame-dragging effect has been evaluated by Lense and Thirring as extremely weak
in comparison with the geodetic effect when the bodies considered are the Sun and
one of its planets.
According to Phillips [205], the Gravity Probe B (GP-B) mission is based on a
satellite, launched on 20 April 2004, which carries four spinning gyroscopes pointed
towards some distant star as a fixed reference point. In the absence of the
frame-dragging effect, the axis of each gyroscope would continue to point to the star. On
the contrary, due to the distorted space-time, the direction of the axis drifts with time.
The geodetic drift rate and the frame-dragging rate measured by the four cryogenic
gyroscopes are, respectively, −6601.8 ± 18.3 mas/year and −37.2 ± 7.2 mas/year,
to be compared with the corresponding predicted values of -6606.1 mas/year and
−39.2 mas/year.
The secular relativistic perturbations on the motion of artificial satellites orbiting
around the Earth have been considered by Cugusi and Proverbio (206). These
authors suppose that a satellite moves under the action of the Newtonian central
force and a perturbing force which represents the relativistic deviation from the
Newtonian law of motion. Then, using the method of variation of osculating ele-
ments and taking account only of the first-order perturbations, they obtain the
following expressions (truncated at the first-order powers in xE):
Da ¼ De ¼ Di ¼ 0
1
6plE 24pl2E xE R2E cos i
Dx ¼ 2 ¼ D1 x þ D2 x
c að 1 e 2 Þ 3
5c2 a2 ð1 e2 Þ2
3
1
8pl2E xE R2E
DX ¼ 3
5c2 a2 ð1 e2 Þ2
3
1
6pa2 l2E
1
h 1 i
Ds ¼ 1 5 2 1 e 2 2
c2 ð 1 e2 Þ 2
where lE (km3/s2), RE (km), and xE (rad/s) are, respectively, the gravitational

parameter, the radius, and the angular velocity of the Earth, the subscript E denotes
the Earth, c (km/s) is the velocity of light, and a (km), e, i (rad), X (rad), x (rad),
and s (s) form the set of orbital parameters of the satellite, with s indicating the time
of perigee passage, and D indicating the relativistic correction to the Newtonian
value of each quantity, per revolution of the satellite (i.e. the secular perturbation).
Cugusi and Proverbio point out that only the last three parameters (i.e. X, x, and s)
are subject to secular relativistic perturbations. However, when the periodic (not
only the secular) relativistic perturbations are taken into account, then only the
major semi-axis remains unperturbed [206].
In the preceding expressions, the periapsis shift consists of two terms. The first
term (D1x) relates to the geodetic (or Einstein’s) effect, and the second term (D2x)
relates to the frame-dragging (or Lense and Thirring’s) effect, which is due to the
rotation of the Earth around its axis. Therefore, D2x does depend on xE, whereas
D1x does not. Because of the presence of the cos i factor in the second term, D2x is
opposite in sign to D1x when the inclination (i) of the orbital plane with respect to
the equator is less than 90°. The relativistic correction DX to the right ascension of
the ascending node also depends on xE, but is not affected by the inclination (i).
The relativistic correction Ds to the time of perigee passage does not show a
dependence on xE in the expressions given above. However, Ds depends actually
on x2E, through a term which has been omitted above.
Cugusi and Proverbio tabulate the values of the relativistic perturbations for the
planets of the Solar System, for the Moon, and for some artificial satellites of the
Earth. They show that the relativistic perturbations relating to artificial satellites are
higher than those for planets. In particular, the geodetic effect (D1x) in the shift of
the perigee for an artificial Earth satellite is larger (up to 30 times) than the same
effect in the shift of the perihelion for Mercury. In addition, the two perturbations
induced in the orbits of artificial satellites by the angular velocity of the Earth (xE)
are about three orders of magnitude higher than those induced in the orbit of
Mercury by the angular velocity of the Sun (xS).
Nevertheless, the relativistic perturbations in the motion of artificial satellites
orbiting around the Earth are very weak (about 10″/year on x, 0″.2/year on X, and
0 s.2/year on s) in comparison with the perturbations arising from the causes
considered in the preceding sections.
Seeber [207] agrees with Cugusi and Proverbio [206] on the comparatively small
contribution given by the relativistic effects to the total perturbations affecting the
orbits of Earth artificial satellites. The corresponding perturbing acceleration vector
(arel) to be taken into account in this case is given by Seeber, on the basis of the
IERS Technical Note No. 21 [44], as follows

lE 4lE
arel ¼ v2 r þ 4ðr vÞv
c2 r 3 r
where r and v are, respectively, the position vector and the velocity vector of the
satellite. According to Seeber, the magnitude of the perturbing acceleration vector
arel due to relativistic effects is about 3 10−10 m/s2 for GPS satellites and
1 10−8 m/s2 for the TOPEX/Poseidon satellite.
3.25 The Perturbations Due to Continuous Low-Thrust

Propulsion
Perturbations may intentionally be induced in the orbit of a spacecraft for the

purpose of changing one or more of its orbital elements. In particular, according to
Biblarz [208], the perturbations generated artificially by means of electric propul-
sion devices can be used to:
3.25 The Perturbations Due to Continuous Low-Thrust Propulsion 549
• overcome translational and rotational perturbations, of natural origin, in satellite

orbits (examples are the north–south station-keeping of satellites placed in
geostationary orbits, performed by means of thrusts acting perpendicularly to the
orbital plane, for counterbalancing the effects of luni-solar perturbations which
would otherwise increase the orbit inclination by about 0.85 degrees/year, and
the east–west station-keeping of such satellites, performed by thrusts acting
tangentially to the orbit, for counterbalancing the effects of the ellipticity of the
Earth equator, which would otherwise cause the satellite to drift, by moving it
from its assigned longitude);
• increase the orbital speeds of satellites, for the purpose of raising their orbits from
a low-Earth orbit (between about 200 to 2000 km in altitude) to higher-altitude
orbits or even to a geostationary orbit (35786 km in altitude); and
• perform interplanetary travels, such as those of deep-space probes.
The impulsive manoeuvres which can be used to that effect will be considered at
length in Chap. 4 of this book. The present section is meant to show the perturbing
effects of low-thrust manoeuvres.
The operating principle of an electric ion thruster is shown in the preceding

figure, due to the courtesy of NASA–JPL [209]. The electric power coming from an
external source (usually solar panels) ionises the propellant (Xenon). The ions so
obtained are accelerated by an electric field between two grids. On leaving such
grids, electrons are injected by a cathode into the beam of accelerated ions, in order
to maintain a neutral plasma. As has been shown by Stuhlinger [210], an electric
propulsion engine produces a low thrust over a long period of time. This is because
its power-to-thrust ratio
1 dm 2
P 2 dt ve 1 1
¼ ¼ ve ¼ g0 Isp
F dm 2 2
ve
dt
depends on the electric power required to accelerate the exhaust stream of pro-
pellant particles which produce the thrust.
In the preceding expression, P stands for power (kinetic energy of jet per unit
time), F for thrust, dm/dt for mass flow rate, ve for the average velocity of the
exhaust stream along the axis of the engine, g0 = 9.81 m/s2 for the acceleration of
gravity at the reference altitude, and Isp for specific impulse.
The expression written above indicates that the power-to-thrust ratio is pro-
portional to the exhaust velocity of the propellant (or to the specific impulse of the
thruster). Therefore, high values of specific impulse imply high values of power per
unit of thrust, and hence high values of electric power supply.
Account must also be taken of the thruster efficiency, η, which is
P
g¼
PT
that is, the jet power generated by the thruster divided by the total electric power
supplied to the thruster itself. Since the value of η is less than unity, a part of the
electric power supplied is converted into jet power, and the other part is wasted.
For example, the power-to-thrust ratio for an electric propulsion engine having a
specific impulse of 2800 s and an efficiency of 65%, as is the case with a Xenon ion
thruster (according to Martinez-Sanchez and Pollard [211]), is
P 1
¼ g0 Isp ¼ 0:5 9:81 2800 ¼ 13734 m/s
F 2
Such an engine would require an electric power
P 13734
PT ¼ ¼ ¼ 21129:23 W
g 0:65
in order to provide a thrust of 1 N. If solar panels generating 0.03 W/cm2 for a

satellite orbiting at a distance of 1 AU from the Sun were available, then two wings
of 2.3 m 15.3 m = 35.2 m2 each would be necessary to the satellite for the sole
purpose of supplying this power. This is a very demanding requirement.
By comparison, the Dawn spacecraft, which has been launched from Cape
Canaveral on 27 September 2007 and uses an ion propulsion engine on its long
journey from the Earth to the asteroid belt, has its solar cells placed on two solar
arrays of 2.3 m 8.3 m = 19.09 m2 each, which provide 10.3 kW at a distance of
1 AU (1.3 kW at end-of-life at a distance of 3 AU) to drive the spacecraft
(22–35 V) and the solar electric ion propulsion system (80–140 V). Its engines
have a specific impulse of 3100 s and a thrust of 90 mN [212].
Consequently, the thrust generated by most electric propulsion engines is less

than 1 N. The following figure, due to the courtesy of NASA–JPL (209), shows the
specific impulse and thrust of different propulsion systems.
The electric propulsion systems are limited in power, because the rate at which the
energy coming from the external source is supplied to the propellant is limited by the
mass of the power system. On the other hand, the low thrust-to-mass ratio peculiar to the
electric propulsion systems makes them apt to be operated for long periods of time, going
from hours to years. They also have exhaust velocities and specific impulses, whose
values are much higher than those of chemical rockets, as shown in the preceding figure.
According to Biblarz [208], the electric propulsion systems may be classified in
the following three fundamental categories:
• electrothermal devices, wherein a propellant is heated electrically and expanded at
supersonic speeds through a nozzle, like the propellant used in a chemical rocket;
• electrostatic devices, wherein acceleration is obtained by the interaction of
electrostatic fields with charged propellant particles, such as atomic ions, dro-
plets, or colloids; and
• electromagnetic devices, wherein acceleration is obtained by the interaction of
electric and magnetic fields within a plasma.
As has been shown above, the accelerations made available by electric propulsion
systems are too low to overcome the gravitational field acting on a spacecraft at the
moment of its launch from the Earth. On the other hand, such systems are apt to be
used in space, where the pressure level (so low as to be near to vacuum) makes it
possible to the exhaust particles to flow. At present, the principal uses of the electric
propulsion systems are the station-keeping and the attitude control of satellites
orbiting around the Earth. However, as has been shown above for the Dawn
spacecraft, these systems have also been used for interplanetary missions.
Several authors (e.g. those indicated in Refs. [208–215]) have considered the
low forces provided by electric propulsion systems as perturbations of a central
gravitational field. Such authors have shown that the orbital manoeuvres required to
change the orbital elements of an Earth satellite can be performed by using not only
conventional high-thrust rockets but also low-thrust propulsion systems. An
example of these near-Earth applications of electric propulsion is the transfer of a
satellite from a low-Earth orbit to a geostationary orbit.
The mathematical theory, which will be shown below, is due to Edelbaum [213]
and has been exposed at length by Chobotov [214]. The acceleration thrust vector
aT aTr ur þ aT/ u/ þ aTh uh
acting on a spacecraft is expressed in a body-fixed system of reference, whose unit

vectors (ur, u/, and uh) point, respectively, in the radial, transverse, and bi-normal
directions, such that ur and u/ lay in the instantaneous plane of motion of the
spacecraft, and uh is directed perpendicularly to this plane along the instantaneous
moment-of-momentum vector h = r v.
The three components of the thrust acceleration vector may also be expressed in
polar, body-fixed co-ordinates aT, a, and b, as follows. Let
12
aT ¼ a2Tr þ a2T/ þ a2Th
be the magnitude of the acceleration thrust vector aT. Let a be the in-plane (or pitch)
angle between the projection of the vector aT onto the instantaneous plane of
motion and the transverse direction, as shown in the following figure.
Let b be the out-of-plane (or yaw) angle between the vector aT and its projection
onto the instantaneous plane of motion. The three components of the thrust
acceleration vector aT result from aT, a, and b as follows
aTr ¼ aT cos b sin a

aT/ ¼ aT cos b cos a
aTh ¼ aT sin b
Inversely, the angles a and b result from aTr, aT/, and aTh as follows
2 3

aTr 6 aTh 7
a ¼ arctan b ¼ arcsin4 12 5
aT/
aTr þ aT/ þ aTh
2 2 2
The in-plane angle a is positive above the transverse direction u/, and the
out-of-plane angle b is positive in the positive direction of the instantaneous
moment-of-momentum vector h = r v.
Edelbaum [213] has studied the problem of the optimal low-trust transfer
between two circular orbits having different radii and inclinations with respect to
the equatorial plane of the Earth. In case of a spiral trajectory used to transfer a
low-thrust propelled satellite from a circular low-Earth orbit to a geostationary
orbit, it is necessary to determine the velocity change, Dv, and the transfer time, t,
necessary for the transfer. Assuming a tangential in-plane component of thrust (that
is, a = 0), it is also necessary to determine the out-of-plane angle b in order for the
satellite to reach simultaneously the zero-degree inclination and the geostationary
radius. The angle b is kept piecewise constant during each orbit, that is, the
out-of-plane component of thrust reverses direction once per orbit (at 90° to the line
of nodes) and has an amplitude increasing with the orbit radius [211]. Following
Chobotov [214], Edelbaum’s analysis gives the following results.
Let v0 and vf be the velocities of the satellite in, respectively, the low-Earth orbit
and the geostationary orbit. For example, assuming a height h0 = 300 km and an
inclination angle i0 = 28°.5 (corresponding to the latitude of Cape Canaveral) for
the initial low-Earth orbit, there results

12
12
lE 398600
v0 ¼ ¼ ¼ 7:7258 km=s
r0 6378:1 þ 300

12
12
lE 398600
vf ¼ ¼ ¼ 3:0747 km=s
rf 6378:1 þ 35786
Let Di = 28°.5 = 0.49742 radians be the inclination change between the initial
orbit (i0 = 28°.5) and the final orbit (if = 0°.0). Then the tangent of the initial
out-of-plane angle (b0) of the thrust acceleration vector results from ([214],
Eq. 14.38, p. 341):

sin 12 pDi
tan b0 ¼ v0 1
vf cos 2 pDi
After introducing the values indicated above for v0, vf, and Di in the preceding
expression, there results b0 = 0.37242 radians = 21°.338.
The total velocity change results from ([214], eq. 14.73, p. 348):
v sin b0
Dv ¼ v0 cos b0 0
tan 12 pDi þ b0
After introducing the values indicated above in the preceding expression, there
results
Dv ¼ 5:9508 km=s ¼ 5950:8 m=s
The total transfer time is computed by means of ([214], eq. 14.22, p. 338):
Dv
t¼
aT
where aT is the magnitude of the thrust acceleration vector. In the example con-
sidered, we suppose to have to do with a satellite of dry mass mpay = 1200 kg
powered by two solar panels, each of which has a surface of 2.3 m 8.3 m (as is
the case with the Dawn spacecraft).
Therefore, the total area of the solar panels is 2 2.3 8.3 = 38.18 m2. If
these solar panels generate 300 W/m2 at a distance of 1 AU from the Earth, then the
input power from the solar panels is
P ¼ 38:18 300 ¼ 11454 W
Supposing that the satellite is propelled by Xenon ion thrusters having a specific
impulse Isp = 2800 s and a total efficiency η = 0.65, the total thrust available is
2gP 2 0:65 11454

F¼ ¼ ¼ 0:54209 N
Isp g0 2800 9:81
The mass of propellant, mprop, results from Tsiolkovsky’s rocket equation (see,
for example, Goebel and Katz [215]):

Dv 5950:8
mprop ¼ mpay exp 1 ¼ 1200 exp 1 ¼ 290:28 kg
Isp g0 2800 9:81
Assuming a constant magnitude
F 0:54209
aT ¼ ¼ ¼ 4:5174 104 m=s2
m 1200
of the thrust acceleration vector, the total transfer time is
Dv 5950:8
t¼ ¼ ¼ 1:3173 107 s ¼ 152:47 days
aT 4:5174 104
Strictly speaking, the mass m of the satellite is not constant, but decreases with
time because of the propellant consumption, as follows
dm F
m0 ¼
dt Isp g0
with m = m0 at t = 0, and consequently the magnitude of the thrust acceleration

vector also varies with time, because of the decreasing mass. A better approxi-
mation than that made above for computing this magnitude and hence the total
transfer time results from taking a constant value m′ = –F/(Ispg0) of the propellant
mass flow rate.
The magnitude of the velocity vector at any time t (from t0 = 0 to tf) is given by
Ref. [214, Eq. 14.67, p. 347]:
1
v ¼ v20 2v0 aT t cos b0 þ a2T t2 2
where aT is the magnitude and b0 is the initial (at time t0 = 0) out-of-plane angle of
the thrust acceleration vector. At any successive time, t, the current out-of-plane
angle, b, of the thrust acceleration vector results from (214, p. 350)
v0 sin b0
tan b ¼
v0 cos b0 aT t
and the current inclination variation, Di, results from [214, p. 350]

2 aT t v0 cos b0 p
Di ¼ arctan þ b0
p v0 sin b0 2
The inclination, i, of the instantaneous orbital plane with respect to the equator is
given by Ref. [214, p. 350]
i ¼ i0 Di
where the plus sign holds when if > i0, and the minus sign holds when if < i0. In the
example considered above, since if = 0° is less than i0 = 28°.5, then i = i0 − Di.
The results given above follow from the requirement of determining the optimal
out-of-plane angle, b, for minimum velocity change, Dv. The variational integral
involves a single Lagrange multiplier, kv, since it involves a single integral con-
straint equation for the velocity change, while maximising the change in orbital
inclination, Di. The control variable is the out-of-plane angle. The necessary con-
dition for a stationary solution results from setting the partial derivative of the
integrand function of the variational integral with respect to the control variable to
zero. The analytical developments of this problem can be found, for example, in
Chobotov [214].
The Edelbaum method described above dates back to 1961. This method applies
to a satellite placed in a nearly circular orbit around the Earth and subject to a
continuous thrust. Subsequent studies (e.g. those of [216, 217]) have considered the
general case of an Earth satellite placed in an elliptic orbit and subject to a dis-
continuous thrust, which acts along arcs of orbit centred around apogee and perigee.
In these studies, which are described below, a simple pre-defined steering pro-
gramme is used together with an averaging technique which takes in consideration
only the secular rate of change of the orbital elements. In particular, the secular rate
of change of each orbit element is considered by holding the other elements (except,
of course, the true anomaly) constant over one revolution.
Pollard [217] has made an analysis of simultaneous eccentricity and inclination
changes, adjusting the argument of perigee and right ascension of the ascending
node, phase adjustment with control of the right ascension of the ascending node,
and the effect of atmospheric drag on major semi-axis and eccentricity. In this
analysis, the thrust acceleration vector, aT, has the same three components (namely
radial aTr, transverse aT/, and bi-normal aTh) which have been described above, and
a magnitude resulting from
12
aT ¼ a2Tr þ a2T/ þ a2Th
Discontinuous thrusting is performed by means of apogee-centred or

perigee-centred burns over an arc of orbit specified by an angle a, which is shown
in the following figure. In this figure, the two arcs MPN and RAS are centred
around, respectively, perigee, P, and apogee, A. The point F indicates the centre of
the gravitational force acting on the satellite. The major and minor semi-axes of the
elliptic orbit are indicated by, respectively, a and b.
The two arcs MPN and RAS are given in units of eccentric anomaly, Æ, where -
a Æ + a (or 180° − a Æ 180° + a).
Progr. aTr(Æ) aT/(Æ)

1 0 aTr/
2 aTr/(e sin Æ)/(1 − e2 cos2Æ)½ aTr/[(1 − e2)/(1 − e2 cos2Æ)]½
3 aTr/(1 − e2)½(sin Æ)/(1 − e cos Æ) aTr/(cos Æ − e)/(1 − e cos Æ)
4 aTr/(cos Æ − e)/(1 − e cos Æ) −aTr/(1 − e2)½(sin Æ)/(1 − e cos Æ)
One of the four steering programmes indicated in the preceding table (from
[217]) for aTr(Æ) and aT/(Æ) is applied to the satellite:
(1) thrust acceleration perpendicular to the orbit radius vector;
(2) thrust acceleration tangent to the orbit path;
(3) thrust acceleration perpendicular to the major axis of the ellipse;
(4) thrust acceleration parallel to the major axis of the ellipse.
For each of these four in-plane steering programmes, the corresponding in-plane
components, aTr(Æ) and aT/(Æ), of the thrust acceleration vector are indicated in
the preceding table, where aTr/ = (a2Tr + a2T/)½ = aT cos |b| designates the in-plane
component of the thrust acceleration vector aT.
As indicated above, discontinuous thrusting is performed by means of
apogee-centred or perigee-centred burns over an arc of orbit specified by an angle a.
The magnitude, aT, of the thrust acceleration vector and the magnitude, |b|, of the
out-of-plane angle of the same vector are held constant during each burn arc.
The out-of-plane component of the thrust acceleration vector is aTh = aT sin |b|.
When the apogee-centred burn and the perigee-centred burn are, both of them,
used in the steering programme chosen, then the sign of b is reversed twice per
revolution at the minor axis crossing, in order to improve the efficiency of the
plane-change manoeuvres. This requires cyclic yaw manoeuvres or alternate firing of
canted thrusters on opposite sides of the satellite [217]. Pollard also notes that using a
slowly varying b, as is the case with the Edelbaum [213] method shown above, yields
little or no benefit for the examples of application considered by Pollard.
In order to apply Pollard’s method to the case of interest, the user chooses one of
the four steering programmes indicates above, and introduces the corresponding
analytical expressions of the thrust acceleration components aTr, aT/, and aTh (which
depend on either b or Æ) in the Lagrange planetary equations in Gaussian form. Then,
these equations are integrated over the burn arc, in order to determine the change of
the orbital elements during one revolution of the satellite around the Earth. By the
way, it is necessary to also take into account the additional perturbing accelerations
(e.g. those due to the non-spherical shape of the Earth) which cause the satellite orbital
elements to vary with time. The results obtained by so doing must be divided by the
instantaneous orbital period, in order to obtain the secular rates of change of the orbital
elements. These secular rates of change must be integrated numerically to obtain the
velocity increment (Dv) and the transfer time. Pollard uses the following expression to
compute the build-up rate of Dv with time (that is, dDv/dt):
dDv 1 2 12
¼ aTr/ þ a2Th ða þ r e sin aÞ
dt p
with r = +1 for apogee burns and r = −1 for perigee burns, the preceding
expression deriving from the burn duration per revolution

12
a3
tburn ¼ 2 ða þ r e sin aÞ
lE
where a is the major semi-axis of the orbit and lE is the gravitational parameter of
the Earth. Among various examples of application of this method, Pollard considers
the circularisation of a geostationary transfer orbit (which, by the way, is a
high-eccentricity elliptical orbit used in order for a satellite to reach the geosta-
tionary radius at apogee). For this orbit, he considers a perigee radius rP = 6563 km
(corresponding to an altitude of 185 km), an apogee radius rA = 42164 km (cor-
responding to an altitude of 35786 km), and an inclination i = 28°.5 (corresponding
to the latitude of Cape Canaveral).
The thrust acceleration magnitude assumed by Pollard is aT = 3 10−7 km/s2,
corresponding to a thruster producing 50 mN per kW of input power and a satellite
having a power-to-mass ratio of 6 W/kg. The first portion of the transfer uses the
steering programme (1) of those indicated in the preceding table, that is, aTr = 0 and
aT/ = aT cos b, with apogee-centred burns along an arc a = 108°, and an
out-of-plane thrust angle b = +40°.4. This portion increases the major semi-axis to
the desired value of 42164 km in 97 days. The following portion of transfer
uses the steering programme (3) of those indicated in the table, that is, continuous
thrust acceleration perpendicular to the major axis of the ellipse with b = ±26°.0.
This portion reduces the eccentricity and the inclination of the orbit to zero, without
modifying the value of the major semi-axis, in 23 days. Pollard computes the value
of b given above by means of the following expressions
A ¼ ði2 i1 Þð3a þ cos a sin aÞ

B ¼ 2 cos x sin a

e2 þ 1 e1 1
C ¼ ln
e2 1 e1 þ 1

A
tanjbj ¼
BC
where the subscripts 1 and 2 designate respectively the initial conditions and the
final conditions.
In order to compute the velocity increment, Dv, required for circularising the
transfer orbit with a simultaneous change of inclination, Pollard introduces the
value of b computed as shown above in the following expression

2aaT Dt lE 12 2a jarcsin e1 arcsin e2 j
Dv ¼ ¼
p a cosjbj 3a þ cos a sin a
The values of a and b used in the first portion of transfer are so chosen as to
minimise the product Dv Dt of the total velocity increment and the total transfer time
for the selected steering programme. In other words, Pollard seeks a favourable
trade-off between propellant mass and transfer time. He has found that, during the first
portion of transfer, the argument of perigee, x, changes because of both the natural
drift and the out-of-plane component of the thrust acceleration vector. In order to get
around this problem, Pollard begins the transfer with x = −13°.5 so as to have
x −1° when the orbit inclination, i, approaches the zero value. This assures a good
geometrical efficiency for inclination changing throughout the mission. The total
velocity increment computed by Pollard is Dv = 2.50 km/s with a transfer time of
97 + 23 = 120 days for this low-thrust transfer, which value is to be compared with
Dv = 1.84 km/s computed by Pollard in case of high thrust (impulsive manoeuvre).
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Chapter 4
Impulsive Orbital Manoeuvres
Unlike natural celestial bodies, the space bodies which are made by man, such as
artificial satellites and spacecraft, are not designed to remain forever in the orbits
into which they were originally inserted at the moment of their launch.
In addition to natural forces which perturb these orbits, there are artificial forces
exerted by propulsion systems carried aboard space vehicles for the purpose of
changing their orbital elements.
The need for this change arises from several reasons. For example, an artificial
satellite revolving around the Earth and meant for communications purposes is
inserted firstly into a Low Earth Orbit (LEO), by means of a separate launch
vehicle, and then into a Geostationary Earth Orbit (GEO), by means of its own
propulsion system. After reaching the geostationary orbit, the same satellite must be
kept in its desired point of station, by using its propulsion system to overcome the
translational and rotational perturbations of natural origin (see Sect. 3.25).
A military satellite meant for the surveillance of a given territory may be made to
change its original orbit in order to track a different target. A satellite which has
reached the end of its life may be made to either re-enter the Earth atmosphere (as is
the case with a low-altitude satellite) or rest in the so-called graveyard orbit, which
is an orbit about 300 km higher than the geostationary orbit, used for the purpose of
removing an old satellite from the highly congested ring placed around the Earth
equator. Finally, an interplanetary spacecraft is usually subject to several orbit
changes or corrections before being inserted into its operational trajectory.
By orbital manoeuvre we mean any intentional change of the orbital elements of
a spacecraft. The manoeuvres obtained by means of chemical rockets are impulsive,
because they result from brief firings of on-board rocket motors for the purpose of
changing instantaneously the magnitude and direction of the spacecraft velocity
vector. By instantaneous change, we mean that the variation of position of a
spacecraft during the firing is assumed to be negligible in comparison with the

https://doi.org/10.1007/978-3-319-62220-0_4
572 4 Impulsive Orbital Manoeuvres
whole orbit. This implies the use of high-thrust rockets, whose burn times are
negligible in comparison with the orbital period of the spacecraft.
In particular, in accordance with Pustynski [1], the following assumptions are
made in the present chapter:
• the central attracting body and the spacecraft are far enough from other celestial
bodies as to be subject in negligible degree to their gravitational forces;
• the gravitational field due to the mass of the spacecraft perturbs to a negligible
extent the gravitational field due to the mass of the central attracting body, that
is, the mass of the former is much smaller than the mass of the latter;
• the gravitational field due to the mass of the central attracting body is spherically
symmetric; and
• the geometric centre of the central attracting body is in the focus of the ellipse
described by the orbiting body, and this centre can be assumed as the origin of
an inertial system of reference.
Strictly speaking, the conditions indicated above are not fulfilled in practice.
This is because, when the central attracting body is the Earth and the orbiting body
is an artificial satellite, then the satellite is also subject to the gravitational attrac-
tions exerted by the Sun, the Moon, and the other planets. In addition, the gravi-
tational field due to the mass of the Earth is not spherically symmetric, nor are
non-gravitational forces (the atmospheric drag for one) fully negligible.
However, in most cases, these conditions can be considered as fulfilled, at least
in the limits of a first approximation, provided that the firings are performed over
short times by means of high-thrust chemical rockets, such as those which will be
described in the following paragraph.
Assuming the validity of the assumptions made above, the following conclu-
sions can be drawn:
• the orbiting body moves in the gravitational field due to the central attracting
body in compliance with Kepler’s laws (that is, the orbit is a conic section
having the attracting body at a focus, the radius vector sweeps equal areas in
equal times, and the square of the orbital period is proportional to the cube of the
major semi-axis);
• the thrust exerted by the rocket motor is the only force which perturbs the orbit
of the attracted body;
• when the rocket motor is fired instantaneously in only one point of the orbit,
then the attracted body returns to this point in its next revolution, but the shape
of the orbit changes as a result of the thrust, because the point of firing belongs
to the old orbit and also to the new orbit (consequently, in order to transfer a
satellite from a circular orbit to another circular orbit of different radius, at least
two firings are necessary);
• in case of a circular orbit, when the rocket motor is fired in the same direction as
the initial velocity vector, then the spacecraft moves to an elliptical orbit whose
perigee is in the point in which the rocket motor was fired and whose apogee is
180° ahead from the firing point (or to a parabolic orbit when the final velocity
is equal to the escape velocity, or to a hyperbolic orbit when the final velocity is
greater than the escape velocity); by contrast, when the rocket motor is fired in
the opposite direction with respect to the initial velocity vector, then the
spacecraft moves to an elliptical orbit whose apogee is the point of firing and
whose perigee is 180° ahead from the firing point (care must be taken that the
perigee height should not decrease to such an extent as to cause the spacecraft to
crash on the surface of the Earth).
The behaviour described above follows from the energy integral

2 1
v ¼l
2
r a
where v and r are, respectively, the magnitude of the velocity and the radius vector
of the spacecraft at some point of its trajectory, l is the gravitational parameter of
the central attracting body, and a (whose values are positive, zero, or negative,
depending on whether the spacecraft moves along, respectively, an ellipse, a
parabola, or a hyperbola) is the major (or transverse) semi-axis of the conic section
described by the spacecraft. The energy integral written as follows
1 2 l l
v ¼
2 r 2a
shows that the first term and the second term on the left-hand side are, respectively,
the kinetic energy and the potential energy per unit mass of the spacecraft, and the
term on the right-hand side is the mechanical energy per unit mass of the spacecraft.
Since −l/(2a) depends only on a, because l is a constant, then to each increase in
the major semi-axis of the orbit, there corresponds an increase in the mechanical
energy per unit mass of the spacecraft.
By the way, the escape velocity, vesc, of a body subject to the gravitational force
of a central attracting body is the speed at which the mechanical energy per unit
mass of the attracted body is zero:
1 2 l
vesc ¼ 0
2 r
In case of an attracted body placed on the surface of the Earth (considered as a

spherically symmetric body) at the equator, the radius and the gravitational
parameter to be considered in the expression given above are r = 6378.137 km and
l = 398600.44 km3/s2, and then the escape velocity of that body from the Earth is
1 1
2l 2 2 398600:44 2
vesc ¼ ¼ ¼ 11:18 km/s
r 6378:137
All impulsive orbital manoeuvres are based on the following fact: the orbit of
a spacecraft revolving about a central attracting body is completely determined
by specifying the position (r) and velocity (v) vectors of the spacecraft at any
point of that orbit. To any change of the velocity vector of the spacecraft at any
point of its orbit, there corresponds a change of that orbit to a new orbit which
complies with the new velocity vector. In other words, any conic-section orbit
travelled by a spacecraft can be changed to another conic-section orbit travelled
by the same spacecraft, by means of a suitable change in the velocity vector of
the spacecraft.
4.2 Engines and Propellants for High-Thrust Rockets
High-thrust propulsion systems for rockets and spacecraft use chemical propellants,
which are burned in the combustion chamber of a rocket to generate thrust.
According Newton’s third law of motion (to every action there is always opposed
an equal reaction), the continuous ejection of a stream of hot gases in one direction
causes a steady motion of the rocket in the opposite direction. A propellant, in turn,
comprises a fuel and an oxidiser. A fuel is a substance, such as liquid hydrogen (H2)
or a mixture of hydrocarbons, which burns in the presence of oxygen to form the
hot gases which are accelerated and then ejected at high speed through a nozzle. An
oxidiser is either liquid oxygen (O2) or a substance releasing oxygen, such as
hydrogen peroxide (H2O2), nitrous oxide (N2O), nitronium perchlorate (NO2ClO4),
or another substance, such as fluorine (F2), having a strong tendency to accept
electrons. The oxidiser combines with the fuel in a proper mixture ratio, as will be
shown below.
Four categories of high-thrust rocket motors may be identified, according to the
physical state of the propellants carried within such rockets. These categories are
solid-propellant motors, liquid-propellant engines, gaseous-propellant engines, and
hybrid-propellant engines, the last of them being those which use propellants stored
in at least two of the three (solid, liquid, and gaseous) physical states of matter. The
following figure, due to the courtesy of NASA-JPL [2], shows a scheme of a
pumped liquid bi-propellant rocket engine.
The necessity of carrying not only a fuel but also an oxidiser within a rocket is
due to the fact that the rocket must be able to operate inside or outside the atmo-
sphere of the Earth. In the rocket motor illustrated below, the magnitude of the
thrust (F) generated depends on the propellant mass flow rate ðm_ dm=dtÞ through
the motor, the magnitude of the exit velocity (ve) of the exhaust gases, and the
pressure (pe) at the exit from the nozzle.
4.2 Engines and Propellants for High-Thrust Rockets 575
The principle of functioning of a converging–diverging nozzle (also known as a

de Laval nozzle, after the name of its inventor, Gustaf de Laval) is illustrated in the
following figure, due to the courtesy of NASA-JPL [2].
As has been shown by Crown [3], a compressible fluid at virtually zero velocity
in the combustion chamber is accelerated through the converging portion of the
nozzle to sonic speed (M = 1) in the throat where, if the converging portion is
properly designed, the flow is uniform and parallel. The fluid is then expanded in
the diverging portion of the nozzle until the desired Mach number M > 1 is reached
in the test section, where the flow is again uniform and parallel.
In other words, the cross section (perpendicular to the gas flow) of the nozzle
immediately downstream of the combustion chamber decreases, in order for the
subsonic (M < 1) flow to increase its speed up to the speed of sound (M = 1).
The sonic flow condition is reached in the throat, where the cross-sectional area
of the nozzle reaches its minimum value. In order for the flow to increase further its
speed in supersonic conditions (M > 1), the exhaust gases must expand downstream
of the throat, and therefore the cross-sectional area of the nozzle must increase. This
is because, in the converging portion of the nozzle, the maximum uniform velocity
which can be reached by the flow in any section is the velocity corresponding to the
local velocity of sound. Further increases of velocity can only be obtained by
expanding the gas in the subsequent diverging portion of the nozzle. Two types
(bell-shaped and conical) of the diverging portion of a nozzle are shown in the
following figure (courtesy of NASA [4]).
A bell-shaped nozzle has the advantages, over a conical nozzle, of better per-
formance and shorter length. The former has a radial-flow section in the initial
divergent region. This fact generates a uniform flow directed along the symmetry
axis at the exit cross section of the nozzle. In addition, the gradual change of the
wall angle with respect to the axis prevents oblique shocks. Some of the methods
(in particular, those due to Prandtl–Busemann, Puckett, and Foelsch) to determine a
nozzle contour have been described by Crown [3]. Rao [5] has applied the calculus
of variations to determine the shape of a nozzle contour leading to the maximum
thrust. For engineering purposes, near-optimum parabolic contours are suitable for
many applications [4].
There is one additional prerequisite for the establishment and perpetuation of
supersonic flow. This is the maintenance of at least the minimum pressure ratio
between the combustion chamber (whose pressure is pcc) and the test section
(whose pressure is p) of the nozzle, that is,
c
c 1 2 c1
pcc ¼ p 1 þ M
2
where M is the Mach number in the test section, and c = cp/cv is the specific heat
ratio of the exhaust gas. Strictly speaking, the preceding equality holds when the
combustion chamber is a reservoir into which the flow may be decelerated isen-
tropically to zero speed. In other words, pcc is the stagnation pressure of the gas.
The additional condition described above is due to the fact that the velocity of
sound is equal to the velocity of propagation of a pressure wave within a medium.

Therefore, when these two conditions are fulfilled throughout the diverging portion
of the nozzle, it is impossible for a pressure disturbance downstream of the nozzle
throat to influence the flow at the throat or upstream of the throat. In other words,
the exhaust gas expansion at supersonic speeds in the diverging portion of the
nozzle prevents pressure waves from propagating backwards [2].
The section ratio (also called expansion ratio) is the cross-sectional area (Ae) of
the nozzle at exit divided by the cross-sectional area (At) of the nozzle at the throat.
In practice, according to Munday et al. [6], for traditional smoothly contoured
convergent–divergent nozzles, designed for isentropic expansion and parallel exit
flow, matching the exit pressure, pe, to the ambient pressure, p0, produces a
supersonic jet which is shock-free. Such a condition, with pressures matched at the
exit, is referred to as the nozzle design condition. Therefore, a convergent–diver-
gent nozzle, be it conical or traditional (that is, bell-shaped), may be characterised
by a design Mach number, Md, which is the Mach number at the nozzle exit
corresponding to the ratio of areas at the nozzle throat, At, and at the nozzle exit, Ae,
given by Ref. [6]:
2 c þ 1
Ae 1 2 c 1 2 c1
¼ 2 1þ Md
At Md cþ1 2
Since only sonic and supersonic flows are of interest, then Md 1. The
operating condition of an exhaust nozzle may be characterised by a jet Mach
number, Mj, which corresponds to the pressure ratio between the total pressure
supplied to the nozzle, ptc, and the ambient static pressure outside the nozzle, p0,
given by Ref. [6]:
c
k 1 2 c1
ptc ¼ p0 1 þ Mj
2
The nozzle is said to be over-expanded when Mj < Md, on design when

Mj = Md, and under-expanded when Mj > Md.
As will be shown below, for optimum performance in terms of thrust, the gas
pressure, pe, at the exit plane of the nozzle should be exactly equal to the pressure,
p0, due to the environment around it. In the vacuum of space (p0 = 0), this is
impossible. The bigger the exit area, Ae, of the nozzle, the closer the rocket gets to
the optimum thrust. However, at some point, the additional thrust gained is not
worth the added mass which is necessary to make the exit area bigger. In conditions
of optimum expansion ratio, the rocket motor produces the maximum thrust.
A nozzle, which ends before the exhaust gas reaches the pressure of the outside
environment, is called an under-expanded nozzle. In these conditions, the rocket
does not get all the thrust available from its motor, because the expansion ratio is
too low, and therefore the exhaust gas pressure is greater than the environmental
pressure. By contrast, a nozzle, whose cross-sectional area at the exit plane is too
large, is called an over-expanded nozzle. This occurs when the expansion ratio is
too high, and the pressure of the exhaust gas is less than the environmental pressure.
Since the cross-sectional area at the exit plane is too large, then the exhaust gas
completely expands before reaching the exit plane of the nozzle.
In the atmosphere, it is difficult to get the optimum expansion ratio, because the
air pressure changes with temperature and altitude. Since a given expansion ratio
results in the optimum expansion only at a specific altitude, the design expansion
ratio of the nozzle must be selected in such a way as to give the best average
performance during powered flight.
With reference to the preceding figure, due to the courtesy of NASA [7], the
thrust imparted to the spacecraft by the gas flow can be expressed as follows
_ e
F ¼ mv
where m_ is the mass flow rate of the exhaust gas, and ve is its velocity at the exit
_ the thrust increases with
plane of the nozzle. Therefore, for a given gas flow rate m,
increasing velocity at exit.
A spacecraft is subject not only to the thrust, but also to the pressure, p0, due to
the environment around it. At high altitudes, p0 can be assumed to be equal to zero.
In these conditions, the net force acting on the gas contained in the thrust chamber
(comprising the combustion chamber and the nozzle) is the sum of the reactions
coming from the walls and of the reaction of the absolute gas pressure at the exit.
These two reaction forces are opposed, as shown in the preceding figure.
Therefore, according to the second principle of dynamics, the net force acting on
the gas must be equal to the momentum flux out of the thrust chamber, which is
Z0 _
W
ptc dA Ae pe ¼ ve
g0
Atc
where Atc and ptc are respectively the area and pressure of the thrust chamber, Ae, ve,
and pe are, respectively, the area, gas velocity, and static absolute pressure at the
exit plane of the nozzle, W_ is the steady spacecraft weight flow rate, W=g _ 0 is the
rate of change of spacecraft mass with time, and g0 is the acceleration of gravity at
sea level.
The integral, that is, the first term on the left-hand side of the preceding equation,
represent the integral of the pressure forces (resultant) acting on the thrust chamber
(combustion chamber and nozzle) and projected onto a plane normal to the axis of
symmetry of the thrust chamber, as shown in the preceding figure. This integral is
just the force, F, which acts on the thrust chamber (and, therefore, on the space-
craft). Consequently, the preceding equation may also be written as follows
W_
F¼ ve þ Ae pe
g0
This equality holds under the following hypotheses:

• the injection flow velocity of the propellants is negligible;
• the flow is constant and does not change with time during the burn;
• the products of combustion are in chemical equilibrium after the burn;
• the flow of gases through the exit plane of the nozzle is one-dimensional, that is,
all the molecules of the ejected gas move on straight lines which are parallel to
the axis of symmetry of the nozzle;
• the flow is compressible, that is, the gas density is subject to significant changes;
• the flow is isentropic, that is, frictionless and adiabatic (without heat loss), in
such a way as to depend only on the cross-sectional areas of the nozzle; and
• the gas flowing in the nozzle is considered as a perfect gas, so that the
well-known equation of state (p/q = RT) involving pressure (p), density (q), and
temperature (T) can be used, where R is the constant of the specific gas (that is,
the universal gas constant R* = 8.314462 J mol−1 K−1 divided by the average
molar mass of the exhaust gases. For example, if the exhaust gas were water,
H2O, then R = R*/18).

In the thrust equation written above, the quantity W=g _ 0 ve is called the
momentum thrust, and the quantity Aepe is called the pressure thrust. The latter term
indicates that some (but not all) of the total pressure forces, ptc, available in the
thrust chamber contribute to the kinetic energy possessed by the exhaust gases
which flow in the nozzle. In other words, a part of the gas pressure generated by the
release of chemical energy is (whereas another part is not) used to increase the gas
momentum. The greater the part of gas pressure used for this purpose, the greater
the efficiency of the nozzle.
Let us assume now that the rocket is fired at an ambient pressure, p0, greater than
zero, that is, at low altitudes. In these conditions, the pressure forces due to the
environment and acting on the outside of the thrust chamber walls have no effect on
the gas inside. However, such pressure forces subtract a part (Aep0) from the total
pressure thrust (Aepe). Since the exhaust gas flows at supersonic velocity through
the exit plane, Ae, then the ambient pressure, p0, cannot gain access to it. Therefore,
the ambient pressure generates a net unbalanced force onto the projected thrust
chamber area, in the direction opposed to the thrust, of magnitude Aep0. This fact is
taken into account by rewriting the equation of thrust as follows
W_
F¼ ve þ Ae ðpe p0 Þ
g0
_ 0 , is due to the mass

Since the rate of change of spacecraft mass with time, W=g
_ of the exhaust gas, then the general equation of thrust may also be
flow rate, m,
written as follows
_ e þ Ae ðpe p0 Þ
F ¼ mv
which is the expression given in Ref. [2]. In case of a conical nozzle, the exit plane
of the nozzle has a divergence angle a measured from the axis of symmetry of the
nozzle. Therefore, not all the gas momentum is in the axial direction, there being a
non-axial component of the exhaust gas velocity. This loss of axial momentum in a
conical nozzle is taken into account by means of a factor, k, defined as follows
1 þ cos a
k¼
2
where k is the ratio between the momentum of the gases in a nozzle with a
semi-aperture angle a and the momentum of an ideal nozzle with all gases flowing
in the axial direction [8]. With this correction, the thrust of a rocket motor having a
conical nozzle is [9]:
F ¼ kmv
Under the seven hypotheses indicated above, it is possible to compute the

exhaust gas velocity, ve, at the exit plane of the nozzle as a function of the physical
quantities of the gas.
The equations used below for this purpose are those of energy, continuity,
momentum, and state.
The energy equation (also known as the Bernoulli equation) states the principle
of conservation of energy. In the hypotheses made above, let us consider an adi-
abatic flow between two cross sections (1 and 2) of a nozzle perpendicular to its
axis of symmetry, as shown in the following figure (due to NASA [10]).
The energy equation for an adiabatic flow between the two cross sections shown
above can be written as follows [8]
1 v22 v21
h1 h2 ¼ cp ðT1 T2 Þ ¼
2 J
where J is the mechanical equivalent of heat (J = 1 if the SI is used, that is, if heat is
measured in joules), h = u + p/q is the enthalpy per unit mass of the fluid, u is the
internal energy per unit mass of the fluid, cp is the specific heat of the fluid at
constant pressure, and T is the temperature of the fluid.
The enthalpy, h, measures the total energy of a thermodynamic system and
includes the internal energy, u, of the gas and the amount of energy required to
make room for it by displacing its environment and establishing its density, q, and
pressure, p.
The energy equation states that the energy available for heat transfer, h1 − h2, is
converted into kinetic energy, ½(v22 − v21), of the flow. In other words, an amount of
energy in form of heat possessed by the gas is used to increase the velocity of the
flow. The term cp(T1 − T2) indicates the decrease in temperature resulting from the
energy conversion. The specific heat, cp, has been assumed constant and depends
on the composition of the gas resulting from the combustion. In particular, when the
cross-section 1 is located at the exit of the combustion chamber (subscript cc), as
shown above, then there results v1 0, T1 = Tcc, v2 v, and T2 T.
Therefore, for an isentropic flow (p/qc = constant), the energy equation can be
written as follows

c1 2
Tcc ¼ T 1 þ M
2
c
c 1 2 c1
pcc ¼ p 1 þ M
2
1
c 1 2 c1
qcc ¼ q 1 þ M
2
where M = v/a is the local Mach number and a = (cRT)½ is the local sonic velocity
of the flow.
The second of the four equations mentioned above states the principle of mass
conservation. Since the flow is constant, then the rate at which a quantity of mass
passes through a cross section of the nozzle must be independent of the position of
this cross section along the axis of symmetry of the nozzle, that is,
qAv ¼ constant
where q, A, and v are, respectively, the local density of the gas, the local
cross-sectional area of the nozzle, and the local velocity of the flow.
Since the mass flow rate through the nozzle must be constant in all sections of
the nozzle, the continuity equation may also be written as follows
qAv ¼ qt At vt
where the subscript t indicates the throat of the nozzle, where the Mach number is
equal to unity. By using the equations written above, it is possible to express the
area ratio between the local cross-sectional area (A) and the throat cross-sectional
area (At, M = 1) of the nozzle as follows [3]:
2 c þ 1
A 1 2 c 1 2 c1
¼ 2 1þ M
At M cþ1 2
Let the subscripts 1 and 2 designate two cross sections of the nozzle at two
arbitrary points along its axis of symmetry. With reference to these sections, the
equation of energy can be written as follows
1 v22 v21
h1 h2 ¼
2 J
The preceding equation, solved for v2, yields

1
v2 ¼ 2J ðh1 h2 Þ þ v21 2
which can also be written as follows [8]:

( " c1 # )12
2c p2 c
v2 ¼ RT1 1 þ v1
2
c1 p1
By choosing the cross-section 2 as the exit plane of the nozzle (subscript e) and
the cross-section 1 as the exit plane of the combustion chamber (subscript cc), the
second addend (v21 v2cc) within curly brackets becomes negligible in comparison
with the first. Therefore, the velocity of the exhaust gas at the exit plane of the
nozzle is expressible as follows
( " c1 #)12
2c pe c
ve ¼ RTcc 1
c1 pcc
This is because the cross section of the combustion chamber is larger than that of
the nozzle throat. Therefore, the flow velocity, vcc, at the exit plane of the com-
bustion chamber can be neglected in the expression given above. In addition, the
temperature, Tcc, of the combustion chamber and of the nozzle inlet differs very
little, in isentropic conditions, from the stagnation temperature.
The maximum theoretical value of the exhaust gas velocity, ve, at the exit plane
of the nozzle is reached when the pressure value, pe, at the same plane is zero
(infinite expansion). Setting pe = 0 in the preceding equation yields
12
2c
ðve ÞMAX ¼ RTcc
c1
The third of the equations mentioned above expresses the principle of conser-
vation of momentum in a steady one-dimensional flow. With reference to the
following figure, let us consider an infinitesimal particle of gas moving along a
streamline, s, in a steady, one-dimensional flow, that is, in a flow whose velocity, v,
is independent of time (dv/dt = 0).
We consider here, of all forces which may act on the gaseous particle, only those
due to pressure imbalances. Therefore, gravitational, magnetic, and viscous forces
are neglected. By applying the second law of dynamics to the motion P of the
infinitesimal particle of gas along the streamline, the sum of the forces, F, in the
stream direction, s, is equal to the mass, m, of the particle multiplied by the rate of
change, dv/dt, of its velocity, that is,
X dv
F¼m ¼ pdA ðp þ dpÞdA
dt
Since the mass, m, of the infinitesimal particle is equal to its density, q, times the
infinitesimal distance, ds, times the infinitesimal area, dA, then there results
m ¼ qdsdA
In addition, since the velocity vector of the particle is always tangent to the
streamline, then the rate of change of velocity with time is expressible as follows

dv dv ds dv
¼ ¼v
dt ds dt ds
By substituting the expressions of mass and acceleration in the equation written


dv
pdA ðp þ dpÞ dA ¼ ðqdsdAÞ v
ds
that is,
dp ¼ qvdv
which is the momentum equation. The minus sign in front of dp is due to the fact
that the gas particle moves along the streamline, s, from a region of high pressure to
a region of low pressure; therefore, when the velocity of the particle increases with s,
then its pressure decreases. In other words, to each increase in the flow velocity of a
particle, there corresponds a decrease in its pressure. The preceding equation, also
known as the Euler equation, is a particular case of the general Navier–Stokes
equations [11]. Since the flow is not only steady but also isentropic (p/qc = constant),
then the following equality
dp dq
¼c
p q
holds. The two equations written above, combined together, yield

dp dp dq dq
vdv ¼ ¼ ¼ a2
q dq q q
where a = (dp/dq)½ is the velocity of sound in the flowing gas.

Since M = v/a is the Mach number, then by multiplying and dividing the
quantity on the left-hand side of the preceding equation by v, there results
dv dq
v2 ¼ a2
v q
that is,
dv dq
M2 ¼
v q
By combining the preceding equation with the equation (qAv = constant) which
expresses the principle of mass conservation, there results
dv dA
1 M2 ¼
v A
For subsonic flow (M < 1), a decrease in the cross-sectional area (dA < 0) of the
nozzle causes the flow velocity to increase (dv > 0). This is what happens in the
convergent portion of the nozzle, downstream of the combustion chamber, until the
sonic condition (M = 1) is reached in the throat. By contrast, for supersonic flow, an
increase in the cross-sectional area (dA > 0) of the nozzle causes the flow velocity
to increase (dv > 0). This is what happens in the divergent portion of the nozzle,
downstream of the throat. The equations written above (energy, continuity,
momentum, and state) make it possible to express the propellant mass flow rate, m, _
and the thrust, F, in terms of the combustion chamber (subscript cc) fluid and the
nozzle throat (subscript t) area, as will be shown below.
The energy equation for the combustion chamber (where vcc 0) is
1
qcc c 1 2 c1
¼ 1þ M
q 2
When M2 is much smaller than unity, then the preceding equation yields
q qcc, and consequently
pcc 1 v2 qv2
¼ 1 þ M2 ¼ 1 þ ¼ 1þ
p 2 2RT 2p
Hence
1
pcc ¼ p þ qv2
2
The mass flow per unit area is
m_
¼ qv
A
Since M = v/a, a = (cRT)½, and Tcc/T = 1 + ½(c − 1)M2, then

2 312
6 cRTcc 7
v ¼ M4
c 1 25
1þ M
2
and therefore
2 3 cþ1
2ðc1Þ
1
m_ pcc c M 6
2 1 7
¼ qv ¼ 1 4 5
A ðRTcc Þ 1 þ
2 c 1 2
M
2
The preceding equation, written for the throat (where A = At and M = 1), yields
1 cþ1
m_ pcc c2 2 2ðc1Þ
¼
At ðRTcc Þ12 c þ 1
that is,
At pcc
m_ ¼ CðcÞ 1
ðRTcc Þ2
where
cþ1
1 2 2ðc1Þ
CðcÞ ¼ c2
cþ1
is the Vandenkerckhove function. By introducing the expressions of m_ n and ve

derived above in the equation of thrust
F ¼ mv
it is possible to rewrite this equation as follows [12]:

( þ 1"
cc1 c1 #)12
2c2 2 pe c
F ¼ At pcc 1 þ A e ð pe p0 Þ
c 1 cþ1 pcc
which expresses the thrust in terms of the combustion chamber fluid and the nozzle
throat area. In other words, when the pressure thrust is zero, which occurs when
pe = p0, then the thrust is directly proportional to the throat area, At, and nearly
directly proportional to the pressure, pcc, within the combustion chamber.
The chemical reaction, which takes place in the combustion chamber when the
fuel combines with the oxidiser, is an oxidation. When the fuel is hydrogen (H2),
the product of the combustion is water (H2O) in its gaseous state, because of the
high temperature generated when hydrogen burns in the presence of oxygen (O2),
according the well-known reaction
2H2 þ O2 ! 2H2 O þ heat
This means that two moles of hydrogen (molecular mass 2 2 1 = 4) react

with one mole of oxygen (molecular mass 2 16 = 32) to form two moles of
water. In other words, 4 g of hydrogen react with 32 g of oxygen to form 36 g of
water. The corresponding stoichiometric mixture ratio (o/f) between the oxidiser
(o) and the fuel (f) is
o 32
¼ ¼8
f 4
In other words, eight grammes of oxygen are necessary to burn each gramme of
hydrogen.
When used for rockets, hydrogen and oxygen are most often stored in their
liquid state of aggregation, at very low temperatures (cryogenic propellants). This is
because they take less space in the liquid state than in the gaseous state at normal
temperature and pressure conditions.
Hydrogen remains liquid without evaporating at atmospheric pressure at a
temperature not higher than −253 °C, and oxygen does the same at a temperature
not higher than −183 °C [13].
Another cryogenic fuel used for space propulsion is liquid methane (CH4),
which requires a storage at a temperature not higher than −162 °C [13]. The related
chemical reaction with liquid oxygen is [14]:
CH4 þ 2O2 ! CO2 þ 2H2 O þ heat
One mole (that is, 12 + 4 1 = 16 g) of methane reacts with two moles

(2 216 = 64 g) of oxygen to form one mole (12 + 2 16 = 44 g) of carbon
dioxide and two moles (2 21 + 2 16 = 36 g) of water. The corresponding
stoichiometric mixture ratio (o/f) between the oxidiser (o) and the fuel (f) is
o 64
¼ ¼4
f 16
A cryogenic oxidiser which can be used either alone or in combination with

liquid oxygen is liquid fluorine (F2), which requires a storage at a temperature not
higher than −188 °C. Fluorine is extremely toxic and reacts violently with most
substances [13].
Cryogenic propellants are difficult to store over long periods of time, due to the
low temperatures required by them. Therefore, they are not used in military rockets,
which must be taken in a state of readiness for launch over times of several months.
In addition, liquid hydrogen requires a storage volume many times greater than that
required by other fuels, because of its very low density (0.07 g/ml), according to
Ref. [14].
On the other hand, the higher performance made available by cryogenic pro-
pellants makes them particularly attractive when constraints of storage and reaction
time are not critical. For example, according to NASA [15], the high-efficiency
engines aboard the Space Shuttle orbiter use liquid hydrogen and oxygen and have
a specific impulse rating of 455 s. The RL-10 engines on the Centaur, the US first
liquid-hydrogen/liquid-oxygen rocket stage, have a specific impulse of 444 s. The
J-2 engines used on the Saturn V second and third stages, and on the second stage
of the Saturn 1B, also burned the liquid-hydrogen/liquid-oxygen combination. They
had specific impulse ratings of 425 s. By comparison, the liquid oxygen/kerosene
combination used in the cluster of five F-1 engines in the Saturn V first stage had
specific impulse ratings of 260 s. In a few words, liquid hydrogen, in comparison
with other fuels (such as kerosene), yields more power per unit volume of storage.
When the fuel is a hydrocarbon (for example, kerosene), the oxidation reaction
may be represented [14] as follows
2C12 H26 þ 25O2 ! 24CO þ 26H2 O þ heat
Two moles of dodecane (molecular mass 2 12 12 + 2 26 = 340) react

with twenty-five moles of oxygen (molecular mass 25 2 16 = 800) to form
twenty-four moles of carbon monoxide and twenty-six moles of water. In other
words, 340 g of dodecane react with 800 g of oxygen to form 672 g of carbon
monoxide and 468 g of water. The corresponding stoichiometric mixture ratio
(o/f) between the oxidiser (o) and the fuel (f) is
o 800
¼ 2:35
f 340
It is to be noted that kerosene is a mixture of hydrocarbons. The chemical

composition of kerosene depends on its source, but it usually consists of about ten
different hydrocarbons, each containing from 10 to 16 carbon atoms per molecule;
the constituents include n-dodecane, alkyl benzenes, and naphthalene and its
derivatives. Kerosene is usually represented by the single compound n-dodecane. In
particular, RP-1 is a special type of highly refined kerosene which conforms to
Military Specification MIL-R-25576 and is used as a fuel for rocket motors [13].
For example, RP-1 is used in the first-stage boosters of the Delta and Atlas-Centaur
rockets. It also powered the first stages of the Saturn 1B and Saturn V [15].
Hydrocarbon-based fuels do not pose the severe constraints of storage which are
peculiar to cryogenic fuels. On the other hand, as has been shown above, the
specific impulse yielded by them is considerably less than that made available by
cryogenic fuels.
Another class of substances used for space propulsion includes the hypergolic
propellants, that is, specific couples of fuels and oxidisers which ignite
spontaneously and violently when one gets in touch with the other, without the need
for any ignition source. Such propellants remain liquid at normal temperatures.
These properties make them particularly suitable for propulsion and hydraulic
power systems carried aboard spacecraft, and in particular for those which are
meant to be used for orbital manoeuvring. On the other hand, hypergolic propel-
lants are extremely toxic and/or corrosive and must be handled with the highest
care.
The most common fuels used in hypergolic propellants are hydrazine (N2H4),
monomethyl hydrazine (CH3(NH)NH2), unsymmetrical dimethyl hydrazine (H2NN
(CH3)2), and Aerozine 50, the last being an equal mixture of hydrazine and
unsymmetrical dimethyl hydrazine.
According to Nufer [16], the oxidisers used in combination with the fuels named
above are usually nitrogen tetroxide (N2O4) and various blends of nitrogen
tetroxide with nitric oxide (NO). Hypergolic propellants are used in the core
liquid-propellant stages of the Titan family of launch vehicles, and on the Delta
launcher. The Space Shuttle orbiter uses hypergolic propellants in its orbital
manoeuvring sub-system for orbital insertion, main orbital manoeuvres, and
de-orbit. It also uses hypergols in its reaction control system for attitude control
[15]. The efficiency (in terms of specific impulse) reached in the Space Shuttle by
combining monomethyl hydrazine with nitrogen tetroxide ranges from 260 to 280 s
in the reaction control system, and to 313 s in the orbital manoeuvring sub-system,
the higher efficiency of the latter being attributed by NASA [15] to higher
expansion ratios in the nozzles and higher pressures in the combustion chambers.
In a liquid-propellant rocket, such as the one illustrated above, the fuel and the
oxidiser are stored in two separate containers and are sent to the combustion
chamber by means of two separate systems of pipes, valves, and pumps.
As the sequel will show, liquid-propellant engines are more complex than
solid-propellant motors, but have several advantages over the latter. This is because
the former offer the possibility of controlling the flow of propellant towards the
combustion chamber, thereby permitting of throttling, stopping, or restarting the
motor.
Solid-propellant motors are the simplest of all those indicated above. They use a
mixture of solid compounds comprising both the fuel and oxidiser combined
together. A scheme of such motors is shown in the following figure, which is due to
the courtesy of NASA [17]. The rocket illustrated below consists of the following
parts: a casing (that is, a tubular body) which is usually made of various materials
(such as special steels, titanium, and fibreglass), a propellant (also called grain)
which fills the casing and is firmly cemented to the inner surface of it, a hollow core
(also called perforation) running through the propellant which performs the function
of combustion chamber, an igniter which heats the propellant up to its combustion
temperature, a payload, a throat, a nozzle which accelerates the hot gases and
permits them to escape from an exit section placed on the back side of the rocket,
and optionally some fins which are useful for the purpose of maintaining the
stability of motion when the rocket travels through the atmosphere.
The solid propellant of a rocket like that described above, when ignited, burns
from the inside out (side burning), that is, towards the sides of the casing, at a much
higher rate than would be the case if the burn proceeded from one end of the rocket
to the other (end burning), as shown in the following figure.
The result, for a side-burning propellant, is a higher thrust. Since the walls of the
casing are protected from the hot gases by the propellant itself and often by an
insulation layer, then it is also possible to use heat-treated alloys or plastics for the
casing [18]. The igniter in most rockets and missiles comprises two main compo-
nents, which are the squib and the main igniter charge. The squib is a cartridge
which contains combustible powders (usually black powders and magnesium) in
contact with an electrical wire of high resistance. When the firing switch is acti-
vated, an electric current passes through the wire, which is heated and ignites the
powders. These, in turn, ignite the main charge. The main charge then explodes the
igniter container and sends a hot flame to the main propulsion charge [19].
The propellant occupies, before being ignited, the greater part of the volume
available. After ignition, the hot gases generated by the combustion flow in the
hollow core which serves as the combustion chamber before being exhausted
through the nozzle.
The perforation may be shaped in various ways (e.g. a tube, a tube and rod, a
double anchor, a star, etc.) for the purpose of controlling the thrust through the rate
and pattern of the burn.
The following figure, which is due to the courtesy of Braeunig [14], shows some
types of perforation and the corresponding graphs of thrust as a function of time.
Depending on whether the pressure in the combustion chamber increases, or
remains nearly constant, or decreases with time, the thrust may be progressive, or
neutral, or regressive.
Solid propellants having a tubular shape (1) produce a progressive thrust, which
increases with time and drops suddenly to zero at burnout. Those having a
tube-and-rod cross section (2) produce a neutral thrust, which remains constant in
time and decreases rapidly to zero at burnout. Those having a five-point-star cross
section (3) also produce a constant thrust, which reduces gradually to zero at
burnout. Those having a cruciform cross section (4) produce a thrust which reduces
gradually to zero. Those having a double-anchor cross section (5) also produce a
decreasing thrust, which drops suddenly to zero at burnout. Those having a den-
dritic cross section (6) produce a strong initial thrust, which decreases suddenly to
an almost constant level.
The following figure, due to the courtesy of the US Government [20], shows a
further type of perforation, whose object is to increase the surface area of the core
channel, with the view of also increasing the area exposed to the burn and therefore
the thrust. As the burn proceeds into the propellant perpendicularly to the ignited
surface, the pointed contour evens out into a circular one. The resulting thrust
reaches its highest value at the beginning of the combustion and decreases gradually
thereafter.
The fuels used in solid propellants include asphalts, waxes, oils, plastics, metals,
rubbers, and resins. The oxidisers for solid propellants come from two general
sources: the organic (the source of nitrocellulose and nitroglycerin) and the inor-
ganic (the source of such chemicals as sodium nitrate and potassium perchlorate).
Usually the oxidiser is crystalline in form (like salt or sugar) and is embedded in the
fuel base [21]. The polyurethane fuel base of the most common solid-fuel mixture is
a type of synthetic rubber, which has about the same consistency as that of tyre
rubber [21]. The mechanical properties possessed by a solid propellant must be so
good as to make it apt to bear the stresses coming from ground handling and flight
environment. A solid propellant must also be free from internal bubbles and surface
cracks, because the exposure of a larger burning surface than the intended one
might cause combustion to take place in the cracks, which fact in turn could
generate an explosion [18].
Solid propellants are stable and easily storable. Therefore, they are used for
military purposes, for example, in the submarine-launched ballistic missiles, the
Minuteman intercontinental ballistic missiles, and the first three stages of the MX
missile [21]. They are also used for space operations. Examples are the final stage
of a launch vehicle and the component of a payload which are used to raise the orbit
of a satellite. The Payload Assist Module (PAM) and the Inertial Upper Stage
(IUS) have been used to place satellites from LEOs into geostationary orbits or on
planetary trajectories. In particular, the PAM-DII has been used for Delta and Space
Shuttle payloads. The IUS has been used on the Space Shuttle and the Titan III and
Titan IV class of launch vehicles [15].
One of the US launch vehicles, the Scout, has used exclusively solid propellants.
This four-stage rocket has been used to launch small satellites to orbit. The Titan,
Delta, and Space Shuttle vehicles have used solid-propellant rockets to provide
added thrust at lift-off [15].
The Space Shuttle has used the largest solid-propellant motors built and flown so
far. Each of the reusable boosters contains 453600 kg of propellant, in the form of a
hard, rubbery substance, whose consistency is like that of the eraser on a pencil
[15]. The solid rocket boosters of the Space Shuttle are illustrated in the following
figure, due to the courtesy of NASA [22].
As shown in the figure, the four central segments of each booster are those which
contain propellant. The uppermost one has a star-shaped, hollow channel along its
axis of symmetry. This channel extends from the top to about two-thirds of the way
down, where it gradually rounds out until the channel takes the shape of a cylinder.
This opening connects to a similar cylindrical hole through the centre of the second
through four segments. The propellant, when ignited, burns on all exposed surfaces,
from the top to the bottom of all four segments. Since the star-shaped channel
provides a larger surface exposed to combustion than the simple cylinder which is
in the lower three segments, the total thrust is greater at lift-off, and gradually
decreases as the points of the star burn away, until that channel also becomes
cylindrical in shape. The propellant in the star-shaped segment is also thicker than
that in the other three [15]. The oxidiser contained in the solid propellant used by
the Space Shuttle is ammonium perchlorate (NH4ClO4), which forms the 69.93% of
the mixture. The fuel is a form of powdered aluminium (16%), with an iron oxidiser
powder (0.07%) as a catalyst. The binder which holds the mixture together is
polybutadiene acrylic acid acrylonitrile (C10H13NO2, 12.04%). In addition, the
mixture contains an epoxy-curing agent (1.96%). The binder and epoxy also burn as
fuel, adding thrust. The specific impulse of the Space Shuttle solid rocket booster
propellant is 242 s at sea level and 268.6 s in vacuo [15].
Just as is the case with liquid-propellant engines, so in solid-propellant motors hot
gases are generated by a chemical reaction between a fuel and an oxidiser. Such
gases are accelerated to supersonic velocities through a nozzle, which is designed in
such a way as to produce the resultant thrust force by the conversion of the thermal
energy produced by the chemical reaction into the kinetic energy of the expanded
exhaust gases. In the simple case of a stationary motor generating a one-dimensional
steady flow of exhaust gases, the thrust force, F, has the same expression as that
given above for liquid propellants, that is,
F ¼ mv
where m_ is the mass flow rate at the exit plane of the nozzle, ve is the gas velocity at
the same plane, Ae is the area of the cross section at the same plane, pe is the gas
pressure at the same plane, and p0 is the ambient pressure. Of course, in case of a
conical nozzle with a semi-aperture angle a, the first addend on the right-hand side
of the preceding equality must be multiplied by k = ½(1 + cos a).
Of the quantities named above, only p0 does not depend on the propellant and
motor characteristics. The exhaust gas velocity at the exit plane of the nozzle
has the same expression as that shown above in case of liquid-propellant engines,
that is,
( " c1 #)12
2c pe c
ve ¼ RTcc 1
c1 pcc
where Tcc and pcc are, respectively, the temperature and pressure of the gases in the
combustion chamber, R is the specific gas constant, and c = cp/cv is the specific heat
ratio of the gases. Kuo and Acharya [9] note that the preceding expression of ve
holds when the combustion products of solid propellants do not contain metal
particles. By contrast, in the event of solid propellants containing powered alu-
minium, there is a significant amount of condensed-phase products, such as alu-
minium oxide. In such cases, ve must be computed by means of an expression
which takes account of the condensed-phase mass. Hill and Peterson [12] give the

n 12
pe
ve ¼ 2 Xcs þ ð1 X Þcp Tcc 1
pcc
where Tcc and pcc are, respectively, the temperature and pressure of the gases in the
combustion chamber, X is the fraction of total mass flow in solid (or liquid) form,
1 − X is the fraction of total mass flow in gaseous form, cs is the specific heat of
solid, cp is the specific heat at constant pressure of gas, pe is the gas pressure at the
exit plane of the nozzle, n = R/{[X/(1 − X)]cs + cp}, and R is the specific gas
constant.
According to Hill and Peterson [12], the rate m_ g of generation of gaseous pro-
pellant is equal to the rate of consumption of solid material, and can be expressed as
follows:
m_ g ¼ Ab rb qp
where Ab is the area of the burning surface of the propellant, rb is the recession rate
of the burning front, and qp is mass density of the propellant.
The surface recession rate, rb, also known as the burning rate, is empirically
determined as a function of the propellant composition and certain conditions
within the combustion chamber. Usually, rb is approximated by means of the
following expression [12, 23]
rb ¼ a pncc
which is also known as de Saint Robert’s law or Vieille’s law, where a and n are
constants which are determined empirically for a particular propellant formulation,
and pcc is the pressure within the combustion chamber. Kuo and Acharya [9] note
that Vieille’s law is valid only when the gas cross-flow velocity over the propellant
surface is so low that shear-induced effects on the burning rate can be neglected.
This is because, for most propellants, certain levels of local gas velocity parallel
to the burning surface lead to an increased burning rate of propellant [23]. In other
words, the fast flow of hot gases along the burning surface of the propellant
increases the burn rate, when the flow velocity exceeds a threshold value. This
phenomenon, known as erosive burning, is particularly strong near the nozzle,
where the gas flow increases its speed.
Therefore, in the case of high-velocity flow of combustion gases over the surface
of the burning propellant, a different equation must be used in order to take account
of the erosive burning effect. This equation has usually the following form
rb ¼ a pncc ð1 þ dÞ
where d is an augmentation coefficient to be evaluated experimentally [23].

The coefficient a depends of the initial propellant temperature, Tp, and may be
approximated as follows [12]:
A
a¼
T1 Tp
where T1 is the empirically determined detonation temperature, whose value is

much greater than that of Tp, and A is an empirically determined constant for a
given propellant.
The pressure exponent, n, is generally taken in the range 0.2 < n < 0.6 [8].
Cornelisse et al. [24] note that rb lies generally between 3 and 250 mm/s.
Montgomery [21] gives the values of about 0.40 in. per second (10.16 mm/s) for a
typical double-base propellant and about 0.22 in. per second (5.588 mm/s) for a
dense polyurethane composite. Braeunig [14] gives the values of a and n for the
solid rocket boosters of the Space Shuttle. Sutton and Biblarz [8, page 479,
Table 12.1] give the burning rates and the pressure exponents for some operational
solid propellants.
The rate, m_ n , at which the exhaust gas flows through the nozzle, in the case of a
non-metallised solid propellant, is equal to [9, 12]:
pcc At
m_ n ¼ CðcÞ 1
ðRTcc Þ2
where
2cðc1
þ1
Þ
1 2
CðcÞ ¼ c 2
cþ1
is the Vandenkerckhove function. The preceding expressions of m_ n and ve can be

introduced in the general equation of thrust
F ¼ m_ n ve þ Ae ðpe p0 Þ
This yields [9]:

( " c1 #)12
2c pe c
F ¼ CðcÞpcc At 1 þ Ae ðpe p0 Þ
c1 pcc
The principle of conservation of mass requires the following balance
dM g
m_ g ¼ þ m_ n
dt
between the gas generation rate (on the left-hand side) and the sum of the rates at
which the gaseous mass stored within the combustion chamber increases and the
exhaust gas flows through the nozzle (on the right-hand side). In the preceding
expression, Mg and dMg/dt are, respectively, the mass and the rate of increase of gas
stored within the combustion chamber. As has been shown above
m_ g ¼ Ab rb qp
and
1 cþ1
pcc At c2 2 2ðc1Þ
m_ n ¼
ðRTcc Þ
1
2 c þ 1
Therefore, the principle of conservation of mass requires
1 cþ1
dM g pcc At c2 2 2ðc1Þ
Ab rb qp ¼ þ
ðRTcc Þ2 c þ 1
1
dt
Substituting rb = apncc and dMg/dt = Vcc(dqcc/dt) in the preceding expression

(where Vcc and qcc are, respectively, the volume and the gas density in the com-
bustion chamber) yields
1 cþ1
dqcc pcc At c2 2 2ðc1Þ
Ab qp apncc ¼ Vcc þ
ðRTcc Þ2 c þ 1
1
dt
The time derivative of density can be expressed as follows
dqcc 1 dpcc
¼
dt RTcc dt
because, for the solid propellants commonly used, Tcc is practically independent of
changes in combustion pressure. Hence
1 cþ1
Vcc dpcc pcc At c2 2 2ðc1Þ
Ab qp apncc ¼ þ
ðRTcc Þ2 c þ 1
1
RTcc dt
that is,
" þ 1 #2
cc1 1
Vcc dpcc c 2
¼ Ab apcc qp qcc At pcc
n
RTcc dt RTcc c þ 1
The preceding equation makes it possible to compute the rate of change, dpcc/dt,
of pressure within the combustion chamber in the transient start-up phase, when the
pressure increases rapidly until the standard operating level is reached.
Subsequently, in the second phase, when the pressure, pcc, within the combus-
tion chamber is constant with time, then the left-hand side of the preceding equation
vanishes, and pcc results from
8 91
>
> >
>
1n
>
> >
>
>
<A >
=
ba qp qcc
pcc ¼
>
> At " þ 1 #2 >
cc1 1
>
>
> c 2 >
>
>
: >
;
RTcc c þ 1
Since qcc is negligible with respect to qp [12], then a(qp − qcc) aqp. In the
steady-state phase, characterised by constant pressure within the combustion
chamber, the burning area, Ab, remains constant.
In the third and last phase, also called tail-down phase, after the consumption of
the solid propellant, the pressure, pcc, within the combustion chamber goes rapidly
but not suddenly to zero, because some fragments of propellants are still present
within the combustion chamber. By setting the burning area, Ab, equal to zero in the
general equation, there results
" þ 1 #2
cc1 1
Vcc dpcc c 2
¼ At pcc
RTcc dt RTcc c þ 1
In the third phase, the pressure, pcc, within the combustion chamber decreases
exponentially to zero, as can be shown by solving analytically the preceding dif-
ferential equation.
Solid-propellant rockets have advantages and disadvantages in comparison with
liquid-propellant rockets. According to Sutton and Biblarz [8], the principal
advantages of the former are the simpleness of design (few or no moving parts),
the easiness and readiness of operation, the absence of spills, leaks, or sloshes, the
storability (for a period going from 5 to 25 years) of their propellants, and the
possibility of being recovered, refurbished, and reused, as is the case with the solid
rocket boosters of the Space Shuttle. On the other hand [8], their principal disad-
vantages are the higher danger of explosion and fire (which events may have
catastrophic consequences), the difficulty of handling, the necessity of an ignition
system, the impossibility of being tested before use, the necessity in most cases of a
layer of insulating material, and the difficulty, in ordinary cases, of changing after
ignition a predetermined level and duration of thrust.
On the last issue, solid-propellant rocket motors are commonly believed to be
incapable of throttling, stopping, and restarting. However, according to Ref. [25],
Aerojet (now owned by GenCorp) developed in 1964 a variable-thrust, stop-start,
solid-propellant rocket motor. This motor used a hydraulically actuated plug with
full 2-in. stroke that produced shutdown in 0.5 s by reducing expansion ratio.
According to Montgomery [21], it is possible to stop thrust in a solid-propellant
motor by injecting a high-pressure inert (or neutral) gas into the combustion
chamber, and then restart a motor stopped in this manner. However, Montgomery
adds that such arrangements have not proven worth the effort. Recently, McCauley
et al. of NASA-MSFC Huntsville, AL, [26] have proposed an innovative throttable
solid-propellant rocket motor to be used in thrusters for Mars surface landing
vehicles. Their proposal is based on a mechanism which can control the nozzle
throat area. The mechanism is a tapered plug, called a pintle. The pintle can be
translated in and out of the throat, so as to cause the annular throat area to decrease
or increase the thrust capability of the rocket motor. The propellants used in this
special type of rocket motors are inert unless ignited by an electric potential of
sufficient magnitude. They can be throttled and extinguished by removing the
electric power and then re-ignited by applying again electric power. The fuel used
in these rockets is (poly)vinyl alcohol of high molecular weight, which swells and
intermingles to create a tough plastisol solid when added to the ionic liquid oxi-
diser. The latter is hydroxyl-ammonium nitrate (NH3OHNO3) stabilised with 5%
ammonium nitrate (NH4NO3 or S-HAN-5), which is conductive and generates heat
in the propellant in the presence of an electric current. Unlike most solid propel-
lants, the plastisol cannot be ignited by spark or flame [26].
These special rocket motors require a control system for pressure control,
acceptable slew rate, and pintle position control. In addition, they require a power
source to provide the voltage needed to ignite the propellant. In case of repeated
operations of the thrusters, a replenishable power source is needed. According to
the authors [26], the ability to ignite and extinguish an electrically controlled
extinguishable solid propellant has been demonstrated using thrusters with a
diameter of 0.125 in. (3.175 mm). However, Sutton and Biblarz [8] observe that the
experience in this field is limited.
The third category of rocket motors described here includes the hybrid-
propellant rockets. In such rockets, one component (usually, the oxidiser) of the
propellant is stored in gaseous or liquid phase, and the other (usually, the fuel) is
stored in solid phase, as shown in the following figure, due to the courtesy of
NASA [27].
The oxidiser (liquid oxygen) used in the rocket illustrated above is stored in its
tank under the pressure exerted by liquid helium. Instead of liquid oxygen, nitrous
oxide (N2O) may also be used, which is an oxidising liquefied gas, stable and inert
at room temperature, non-flammable, and deflagrating at a temperature around
600 °C at a pressure of 21 atmospheres. It is also non-toxic, easy to store, and
self-pressurising (its vapour pressure is about 50.8 atmospheres at 20 °C). Further
oxidisers than those named above are gaseous oxygen (O2), hydrogen peroxide
(H2O2), dinitrogen tetroxide (N2O4), and nitric acid (HNO3). In case of nitrous
oxide being used as the oxidiser, there is no necessity of pumps to raise the pressure
within the combustion chamber to the desired value of about 37 atm. When other
oxidisers than nitrous oxide are used, the pressurisation may be provided, in case of
need, by either an inert gas or a turbo-pump system [8].
The oxidiser coming from its tank is injected into a pre-combustion or vapori-
sation chamber located upstream of the solid fuel. The latter, as shown in the figure,
contains several combustion perforations (also called ports), within which fuel
vapour is generated. To this end, a source of ignition, such as a pyrotechnic igniter,
is applied to the fuel grain. Of course, an igniter is needed when self-igniting
(hypergolic) propellants are not used. When the propellant is hypergolic, it is
sufficient to spray the liquid oxidiser onto the solid fuel to initiate ignition. Since no
separate igniter or firing mechanism is needed, then a hypergolic-propellant hybrid
motor could be switched on and off and modulated in thrust by simply operating a
single valve controlling the oxidiser flow. However, Jain [28] points out to this
regard that most often the solid fuels in cast form do not ignite spontaneously with
the usual liquid oxidisers, and that, even when they do, there is a long ignition
delay, which leads to a hard start of the engine. At present, hypergolic hybrid
systems using common storable oxidisers, having short ignition delays, and pro-
ducing innocuous combustion products are rare. Consequently, all recent hybrid
systems use non-hypergolic propellants [28]. An example of a hybrid rocket, which
uses (cryogenic) liquid oxygen (LOX) as its oxidiser and Hydroxyl-Terminated
Poly-Butadiene (HTPB) with added poly-cyclopentadiene (PCPD) as its fuel, is
given by Sutton and Biblarz [8]. The PCPD is added to the HTPB for the purpose of
increasing the fuel density by about 10% over that of the HTPB alone. A mixture of
Tri-Ethyl Aluminium (TEA) and Tri-Ethyl Borane (TEB) is injected in the
pre-combustion (or vaporisation) chamber of the motor indicated above. This
mixture ignites spontaneously on contact with air in the combustion chamber,
vaporising fuel in the dome region. Subsequent injection of liquid oxygen com-
pletes the ignition of the motor [8].
Fuels commonly used with the oxidisers indicated above are Polyethylene (PE),
Poly-Methyl Methacrylate (PMMA or Plexiglas), Poly-Vinyl Chloride (PVC), and
Hydroxyl-Terminated Poly-Butadiene (HTPB). The combination of nitrous oxide
as the liquid oxidiser with HTPB rubber as the solid fuel has been used in the
vehicle SpaceShipOne, which won the Ansari X Prize [13]. According to Sutton
and Biblarz [8], the propellant system of choice for large hybrid booster applica-
tions is liquid-oxygen (LOX) oxidiser combined with Hydroxyl-Terminated
Poly-Butadiene (HTPB).
The vaporised fuel reacts with the oxidiser injected into the combustion chamber
and burns along the length of the fuel grain. The head end and the aft end attached
to each side of the main fuel grain represent, respectively, a pre-combustion
chamber for heating and vaporising the oxidiser, and a mixing chamber for com-
pleting the reaction of the unburned fuel with the oxidiser [29]. The oxidiser flow
can be increased or decreased for throttling. In addition, when the igniter is reusable
and carried onboard, then the motor can be stopped and restarted during the flight.
Therefore, hybrid-propulsion rockets are well suited to all cases in which their
motors must be throttled, shut down and restarted. Examples of such cases are
missions consisting in putting satellites in LEO, increasing the radius of a satellite
orbit from a near Earth to the geostationary value, and also long-duration missions
requiring storable non-toxic propellants [8]. The regression rate of existing solid
fuels used in hybrid-propulsion rockets is lower (about 1/10) than that of solid
propellants [28]. Consequently, high-thrust hybrid rockets have multiple perfora-
tions (combustion ports) in the fuel grain, in order to increase the area of the
burning surface. A wagon-wheel-shaped cross section is shown in the following
figure, due to the courtesy of NASA [30].
The number of combustion ports depends on the desired level of thrust,

acceptable shifts in mixture ratio during burn, motor length and diameter con-
straints, and desired mass velocity of the oxidiser [8]. The fuel regression rate
depends primarily on the oxidiser mass velocity, also called oxidiser flux. The latter
is equal to the mass flow rate of oxidiser in a combustion port divided by the port
cross-sectional area. Therefore, the fuel regression rate depends on the mass flow
rate of oxidiser. Since the fuel grain of a hybrid rocket contains no oxidiser, then the
combustion process and hence the regression of the fuel surface differs from that of
a solid-propellant motor. The mechanisms of heat transfer to the surface of a fuel
grain in a hybrid rocket are convection and radiation. Sutton and Biblarz [8] give
the following expression, to be used in engineering evaluations, for the fuel surface
regression rate
rb ¼ aGn0
where G0 is the oxidiser mass velocity (that is, the oxidiser flow rate divided by the
combustion port area), a is a coefficient determined empirically, and n is an
exponent also determined empirically. Sutton and Biblarz note that the value of rb
varies between 1.27 and 5.08 mm/s, and the value of n varies between 0.4 and 0.7.
They also note that the surface regression rate of some fuels shows little or no
dependence on the pressure within the combustion chamber, whereas the surface
regression rate of other fuels (in particular, of the metallised ones) shows a strong
dependence on the pressure. In particular, the surface regression rate, rb (mm/s), for
the HTPB fuel with oxygen as the oxidiser has been found to be
ðrb ÞHTPB ¼ 0:01149G0:681

0
(Reference [8]). This result has been found for a 50.8-mm-diameter motor having a
10.9-mm combustion port. On the other hand, by using the same fuel and oxidiser,
and for combustion port diameters ranging from 76.2 and 152.4 mm, the surface
regression rate, rb (mm/s), has been found to depend on the combustion port
diameter, Dp (mm), as follows
ðrb ÞHTPB ¼ 0:0002483G0:77

0 Dp
0:71
(Reference [8]). The thrust force, F, of a hybrid rocket has the same expression as
that given above, that is,
F ¼ m_ n ve þ Ae ðpe p0 Þ
where m_ n is the rate at which the exhaust gas flows through the nozzle, ve is the gas
velocity at the exit plane of the nozzle, Ae is the area of the cross section at the same
plane, pe is the gas pressure at the same plane, and p0 is the ambient pressure. The
gas flow rate through the nozzle results from
m_ n ¼ m_ o þ m_ f
where m_ o and m_ f are the flow rates of, respectively, the oxidiser and the fuel.
In order to change the thrust level of a hybrid rocket, it is necessary to change the
flow rate of the oxidiser, by means of the throttling valve placed upstream of the
combustion chamber.
According to Sutton and Biblarz [8], in the simple case of a fuel grain having
only one circular combustion port of radius R, the fuel surface regression rate given
above (rb = aGn0) can be rewritten as follows
n
m_ o
rb ¼ a
pR2
On the other hand, the flow rate of the fuel results from
m_ f ¼ qf Ab rb ¼ qf ð2p RLÞrb
where Ab = 2pRL is the area of the burning surface, and L is the length of the
combustion port. By using the value of rb given above, the preceding equation
becomes
m_ f ¼ 2p1n qf L a m_ no R12n
In the particular case of n = ½, the fuel flow rate does not depend on the radius,
R, of the combustion port and varies as the square root of the oxidiser fuel rate.
When the oxidiser flow rate is constant, then the fuel flow rate increases with
increasing combustion port radius for n < ½ and decreases with increasing com-
bustion port radius for n > ½.
In the general case of a fuel grain having a number N of circular combustion
ports, Sutton and Biblarz [8] also give the values, as a function of time, of the
combustion port radius, the fuel flow rate, the mixture ratio, and the total fuel
consumed.
The hybrid rockets have the advantages peculiar to both the solid-propellant and
the liquid-propellant engines. According to Jain [28], such advantages are the
following:
• safety from the danger of accidental explosions, the fuel and the oxidiser being
stored in different tanks;
• thrust variation control and start-stop-and-restart capability, made available by
the throttling valve;
• insensitivity to possible cracks or defects in the fuel grain, the regression of the
latter being normal to the flow of the oxidiser and not to the surface of the grain;
• higher values of specific impulse than those available with solid-propellant
rockets, and higher values of density impulse than those available with
liquid-propellant rockets;
• insensitivity to ageing, the fuel and the oxidiser coming in mutual contact only
at the moment of their combustion; and
• low cost, well under that of either the solid-propellant or the liquid-propellant
engines.
On the other hand, their disadvantages are:
• lower values (about 1/10) of regression rate than those typical of the
solid-propellant motors;
• slightly lower values (93–95%) of combustion efficiency than those of the other
types of engine; and
• lower values of density impulse than those of the solid-propellant motors, with
the consequent need to occupy larger volumes.
The last section of the present paragraph concerns the thrust vector control of a
rocket or spacecraft. By thrust vector control we mean the ability of a space vehicle
to deflect the direction of its thrust away (that is, at an angle, h, other than zero)
from the longitudinal axis of the vehicle.
The necessity of this control arises from several causes. First, an intentional
change of the direction of the flight path followed by the centre of mass of the
vehicle. Second, an intentional change of attitude (or rotation) of the vehicle around
the principal axes passing through its centre of mass. Third, a correction of either
the flight path or the attitude which becomes necessary in order for the vehicle to
maintain the desired trajectory and orientation. Fourth, a correction of the
misalignment between the thrust vector and the gravity force vector. As to the last
issue, the thrust vector is applied to the nozzle of the vehicle, whereas the gravity
force vector is applied to its centre of gravity. In order for the vehicle not to be
subject to unwanted torques, the directions of these two vectors must be aligned at
180°.
A rocket is also subject when flying through the atmosphere to aerodynamic
forces (lift and drag), which are applied to its centre of pressure, not to its centre of
gravity. These forces produce moments about the principal axes of the vehicle.
These moments, in turn, cause the vehicle to rotate around its centre of gravity.
The centre of gravity, the rotations, and the principal axes of a rocket are shown
in the following figure, which is due to the courtesy of NASA [31]. For a stable
flight, the centre of gravity of a rocket must be above its centre of pressure. The axis
of symmetry of the rocket, or the line joining the tip of the nose with the centre of
the nozzle exit section, is called the roll axis, and a motion of the rocket about the
roll axis is called a rolling motion. The centre of gravity of the rocket lies along the
roll axis. The pitch and yaw axes are mutually perpendicular and form a plane
passing through the centre of gravity and perpendicular to the roll axis. Pitch
moments tend to either lower (if positive) or raise (if negative) the nose of the
rocket. Yaw moments cause the nose to move from side to side.
Both mechanical and aerodynamic methods can be used to redirect the rocket
thrust and provide the necessary steering forces. As long as a rocket travels in the
atmosphere, static fins at the tail of the rocket can generate aerodynamic forces
which act at the centre of pressure and generate moments about the centre of
gravity. These moments oppose the deflection of the rocket in the directions of the
pitch and yaw axes. Of course, the static fins must be so sized as to generate the
amount of aerodynamic forces required to counteract the deviations. On the other
hand, the fins generate further drag, in the direction opposed to the rocket velocity
vector. In spite of that, aerodynamic fins, be they fixed or movable, are very
effective for controlling a vehicle flying through the atmosphere. They continue to
be used in weather rockets, anti-aircraft missiles, and air-to-surface missiles [8].
Generally speaking, a control system used in a rocket may be either passive or
active. A passive control system is a fixed device meant to stabilise a rocket by its
very presence on the outside of the rocket. An example of such a device is a cluster
of fins mounted around the lower end of a rocket, near the nozzle. The purpose of
the fins is to keep the centre of pressure below the centre of gravity of the rocket.
In order to overcome the disadvantages (higher drag and mass) of the fins, active
control systems have been developed. Such systems include jet vanes, movable fins,
jetevators, canards, gimballed nozzles, Vernier rockets, fuel injectors, and
attitude-control rockets. They are illustrated in the following figure, due to the
Jet vanes and movable fins are planar surfaces, which are used in the jet stream
of a rocket to deflect the exhaust gases. Thrust vanes have been employed on the
V-2 and Redstone missiles to deflect the exhaust gas jet of the main propulsion
motors by carbon vanes [33].
Jetevators are control surfaces which can be moved into or against the jet stream
of a rocket, in order to change the direction of the jet flow. The jetevator invented
by Dr. Willy Fiedler was a solid ring with a spherical inside surface which was
hinged over the rocket nozzle [34, page 52]. This device, which was used in the A1
version of the Polaris missile and then replaced by rotatable nozzles in the A2
version, deflects the flow when turned into the exhaust stream. It has the advantage
of not causing propulsion losses when in the neutral position, since, unlike vanes, it
does not interfere with the exhaust flow [33]. Another example of the same device
has been given by Edwards and Parker [35], who describe a semi-spherical shell
hinged to a rocket nozzle and rotated, at the command of a sensing unit, into the
exhaust flow to produce a control force. Jetevators have been most frequently
applied to control the thrust of solid-propellant rockets [33].
Canards are fore-plane surfaces mounted on the front end of a rocket or aircraft.
The following figure, due to the courtesy of NASA [36], shows such surfaces.
Movable fins and canards are quite similar to each other in appearance. The only
real difference is their location on the rocket, because canards are mounted on the
front end, whereas the movable fins are at rear end. In flight, both of them tilt as
rudders to deflect the air flow and cause the rocket to correct its course, in the event
of unwanted directional changes being detected by motion sensors [17].
Another method for changing the direction of the exhaust gases, and hence the
direction of the rocket, is to gimbal the nozzle. In a gimballed thrust control system,
the nozzle of the rocket can be swivelled from side to side. As the nozzle direction
is deflected away from the axis of symmetry of the rocket, so does the direction of
the thrust change with respect to the centre of gravity. This system has been used in
the Space Shuttle and in the Saturn V rocket [17].
In a Vernier thrust control system, there are small rockets mounted outside of the
large thruster. In case of need, such rockets are fired in the proper direction, in order
to produce the desired course change [17].
Attitude-control rockets are used to trim misalignments of the translational
thruster, the basic attitude control with respect to the body principal axes being
provided by a separate set of thrusters. By so doing, the main engine is used only
for the motion of the centre of mass. Attitude-control rockets are used in the
coasting phase of flight. For this purpose, clusters of small rockets are mounted all
around a space vehicle. By firing these small rockets in the proper combination, it is
possible to turn a space vehicle in the desired direction. After the vehicle has been
aimed as required, the main engine is fired, in order to send the rocket in the new
direction. For example, the ascent stage of the Apollo Lunar Module used a fixed,
high-thrust engine for translation and a series of small liquid-propellant rockets for
attitude control, as shown in the following figure, which is due to the courtesy of
NASA [37].
The methods described above use mechanical means to control the thrust vector.
Such means are either aerodynamic surfaces placed within the jet stream of the
main flow or adjustable nozzles used for the purpose of deviating the direction of
the main thrust vector. These methods require actuating components which must
work efficiently in the high-temperature environment of the rocket exhaust and are
invariably associated with a loss of axial thrust when performing thrust vector
control [38].
Other methods use fluidic means to control the thrust vector. These methods use
a static nozzle and the injection (or the removal) of a secondary flow into (or from)
the region between the primary flow, which generates the thrust, and the nozzle.
They do not require any kinematic structure and mechanical actuators. There is a
variety of methods based on fluidic thrust vectoring, which differ one from the other
by the way in which the secondary flow is used for thrust vector control. They may
be classified as follows.
• Co-flow for fluidic thrust vectoring, which is based on the Coanda effect. The
Coanda effect (so named after the Romanian engineer Henri-Marie Coanda) is
the phenomenon in which a stream of fluid, ejected at high speed from a slot and
coming in contact with a convexly curved surface, adheres to that surface along
its curvature rather than continue to travel in a straight line. In other words, the
stream is deflected from the axis of flow and follows the slope (or curvature) of a
divergent wall while increasing in velocity and in mass by entraining additional
fluid [39]. This happens because of the increase in the flow velocity over the
curved surface, which causes pressure to decrease. The low pressure causes not
only the injected co-flow but also the primary flow to be deflected off the nozzle
axis towards the divergent wall, as shown in the following figure.
In the preceding figure, the black arrows represent the primary flow, and the red
arrows represent the injected co-flow. Of all the fluidic methods for controlling
the thrust vector, the co-flow method has been found [40] to be the one which
produces the smallest deflection angle.
• Counter-flow for fluidic thrust vectoring (illustrated below), which is also based
on the Coanda effect to deflect the thrust vector. Strykowski et al. [41] have
shown that the thrust due to a jet stream can be continuously deflected to at least
16° by creating a secondary counter-flowing stream between the primary jet and
an adjacent curved surface. For this purpose, suction is applied asymmetrically
between the trailing edge of a primary nozzle and an aft suction collar, as shown
in the following figure.
This creates a low-pressure region along the suction collar and causes the pri-
mary jet to turn [42]. The results found by Strykowski et al. [41] at Mach 2 show
that the thrust loss is less than 4% and the required mass flow rates are less than
approximately 2% of the primary jet. A co-flow, due to the viscous entrainment

generated by the primary flow, takes place at the wall of the suction collar and
interferes with the primary flow. By activating asymmetrically this co-flow, the
primary flow is deflected towards the side on which the suction flow is applied,
due to the pressure drop which causes the thrust to turn [40].
• Throat shifting for fluidic thrust vectoring, which is performed by injecting
secondary flow at or just upstream of the throat, as shown below.
By so doing, there is no formation of shock waves. This injection skews the

sonic plane and deflects the flow. Another version of this method uses variable
recessed cavities, which deflect the primary flow by means of vortices in the
cavities, as shown in the following figure.
The recessed cavity portion of the nozzle shown above is located between the
upstream minimum area and the downstream minimum area of the nozzle itself.
The fluidic injection is performed at the upstream minimum area. A simulation
study carried out by Deere et al. [43] has shown that substantial thrust vector
angles are obtainable without large penalties in thrust efficiency. This version is
a combined method, because the recessed cavity technique is used in addition to
the throat-shifting technique in order to obtain greater performance.
• Shock vector control for fluidic thrust vectoring, which is performed by injecting
secondary flow downstream of the throat, as shown in the following figure.
The secondary injected flow acts as a compression ramp in the direction of the
primary supersonic flow. This compression induces an oblique shock wave in
the diverging portion of the nozzle at some angle with respect to the direction of
the primary flow. The primary flow, when interacting with the oblique shock
wave, turns away from the axis of symmetry of the nozzle. This turn changes the
direction of the thrust vector. By so doing, the direction of the primary flow does
not change in the vicinity of the throat. The shock vector control may also be
combined with the throat shifting, by using two ports instead of one to inject the
secondary flow into the nozzle, as shown in the following figure.
The aerospace vehicles which use fluidic instead of mechanical thrust vectoring
have the advantages of being lightweight, free from moving parts, less expen-
sive, easy to integrate, and less detectable by radars, the last advantage being of
particular interest in military applications. On the other hand, they must be
designed from the outset as such, and existing vehicles currently in use cannot
be retrofitted. In addition, their capability of directional change is often held to
be lower than that of vehicles using mechanical thrust vectoring. On the last
issue, Strykowski et al. [41] have obtained a thrust vector angle of 16° by using
the counter-flow technique described above, and Wing and Giuliano [44] have
obtained a value of up to 18° for this angle, by using the shock vector control
technique.
Such values are by no means small in comparison with those cited by Sutton and
Biblarz [8, page 611, Table 16.1] for mechanical thrust vectoring systems. They are
lower only than the value of 20°, relating to the movable nozzle (rotary ball with
gas seal). Another disadvantage of the fluidic thrust vectoring methods has been
found in their need for a source of secondary flow. In the event of that source being
the same as that of the primary flow which generates the thrust, it has been argued,
the performance of the vehicle would decrease, at least in the phase of thrust vector
control. As to the loss in performance, Strykowski et al. [41] have found, at Mach 2,
a thrust loss lower than 4%, with the required mass flow rates being lower than
approximately 2% of the primary jet. Consequently, the only strong disadvantage of
the fluidic thrust vectoring methods seems to us to be the deficiency or insufficiency
of data on their behaviour gathered by testing them in flight.
4.3 Launch Windows
The mission of an artificial satellite to be placed in an orbit around the Earth

includes generally three phases. The first of them is the boost phase, during which
the satellite is launched by a rocket from some place on the surface of the Earth and
then inserted into the desired orbit. The second is the spatial phase, during which
the satellite is expected to reach the objectives of its mission. The third is the
re-entry phase, during which the satellite is expected to return safely to the Earth at
the end of its mission.
The present paragraph and the following ones are meant to describe these
phases.
During the boost phase, a rocket is needed for the purposes of overcoming the
gravitational force which keeps the satellite on the surface of the Earth and
launching it out of the atmosphere, at the altitude and velocity which are required
by its orbit. The launch and boost trajectory must satisfy some requirements, which
concern the launch window, the range safety, and the aerodynamic loading. In
addition, the launch and boost trajectory must insert the satellite into its orbit in
such a manner as to place the maximum payload at the desired altitude and velocity
with the minimum weight of the launch vehicle.
The present paragraph concerns the launch window, which is the period of time
during which a satellite must be launched from a given site of the Earth in order to
achieve a desired orbit. This period of time may be expressed either within a single
day or within several successive days, on each of which the launch is possible.
There are three categories of reasons which limit a launch window. The first
includes requirements which depend only on lighting conditions and abort events.
Such is the case with communication, meteorological, or surveillance satellites,
whose orbital planes do not require any specific inertial orientation. By contrast, a
rendezvous mission requires a specific inertial orientation of the orbital plane,
because, in order for the launched vehicle to reach the orbital position of the target
vehicle, the launch must occur around the time in which the orbital plane of the
target vehicle passes through the launch site. This is the second category of reasons
and generates requirements depending not only on lighting conditions and abort
events, but also on rendezvous needs. The third category generates requirements
depending not only on lighting conditions and abort events, but also on a specific
velocity vector, which the launched spacecraft must have with respect to an inertial
system of reference in order to accomplish its mission. Such is the case with a
spacecraft which is meant to reach the Moon or one of the planets. For each of the
three categories indicated above, the related launch window is determined as will be
shown below.
The launch-window requirements generated by the missions belonging to the
first category depend primarily on the necessity of operating in conditions of suf-
ficient daylight at the moment of launch and during the successive period of time, in
the event of an abort during the boost phase, in which case it would be necessary to
reach the emergency recovery site with enough daylight to perform location and
recovery. In order to show such requirements, it is necessary to consider the ori-
entation of the satellite orbit with respect to the rotating Earth.
The inclination, i, of this orbit with respect to the equatorial plane depends on the
latitude ðu L Þ of the launch site (L) and on the azimuth angle (AL) of launch.
According to Robertson [45], by inspection of the following figure and using a
formula of spherical trigonometry, there results
cos i ¼ cos u L sin AL
where u*L i < 90° for prograde orbits.

4.3 Launch Windows 613
As is well known, the condition i u*L is due to the necessity of using as

much as possible the linear velocity
vL ¼ xE RE cos u L
which a spacecraft possesses with respect to an inertial system of reference at the

moment of launch. The spacecraft has this linear velocity because it is launched
from the point L of geocentric latitude u*L located on the surface of a spherical
Earth of radius RE which rotates about the polar axis with angular velocity xE. This
linear velocity is maximum (vL = xERE) when the launch site, L, is at any point of
the equator (u*L = 0°) and is zero when the launch site is at one of the poles
(u*L = 90°). The equality i = u*L holds in case of a launch due East. Prograde
(i < 90°) orbits having inclinations 0 i < u*L require change-of-plane
manoeuvres, to be performed after the spacecraft has been inserted into a parking
orbit of inclination i u*L.
As shown in the preceding figure, u*L denotes the geocentric (not the geodetic)
latitude of the launch site (L). The difference between the two latitudes has been
shown at length in Sect. 2.3. The geodetic latitude (uL) of the launch site results
from a map. For a launch site at the sea level, the relationship between the two
latitudes is the following

tan u L ¼ ð1 f Þ2 tan uL ¼ 1 e2E tan uL
where f = (aE − bE)/aE = 1/298.257223563 [46] and aE = 6378.137 km [46] are,

respectively, the flattening and the equatorial radius of the Earth ellipsoid. The other
quantities, whose approximate values result from the reference values of f and aE
given above, are the eccentricity eE = 2f − f 2 0.08181919 and the polar radius
bE = aE(1 − f) 6356.752 of the Earth ellipsoid. For a launch site having an
altitude other than zero above the sea level, the relationship between the two
latitudes is given by the general formulae given in Sect. 2.3. However, for practical
purposes, the increase in precision (a small fraction of a degree) is negligible.
Likewise, the argument of latitude, uL, of the launch site results from
tan u L
tan uL ¼
cos AL
where for launch sites in the northern hemisphere of the Earth
0\uL \90
ðfor northerly launchesÞ
0\uL \180
ðfor southerly launchesÞ
Let DkL be the projection of uL onto the Earth equator. There results
sin DkL ¼ sin uL sin AL

cos AL
cos D kL ¼
sin i
sin DkL
DkL ¼ arctan
cos DkL
Let Du be the total range, measured along the arc LE of trajectory travelled by
the rocket, to be considered in case of abort. This total range goes from the launch
site (L) to the emergency recovery site (E). The projection, DkE, of Du onto the
equator results from
DkE ¼ arctan½cos i tanðuL þ DuÞ DkL
where 0° < arctan[cos i tan(uL + Du)] < 180° (for northern hemisphere recovery sites)
180° < arctan[cos i tan(uL + Du)] < 360° (for southern hemisphere recovery sites)
The Local Mean Time at a particular meridian (e.g. at the meridian of the launch
site) is defined as 12 h plus the hour angle (the angle measured west of the par-
ticular meridian) of a fictitious Sun that has the same period as the true Sun, but
moves with constant speed along the equatorial plane [47]. The Local Mean Time,
LMT, is related to the Co-ordinated Universal Time, UTC, by
k
LMT ¼ UTC
15
where k is the local longitude, and the plus sign takes effect in case of east lon-
gitude, whereas the minus sign takes effect in case of west longitude. The value of
the local longitude must be divided by 15, because LMT and UTC are expressed in
hours, minutes, and seconds, whereas k is expressed in degrees, minutes, and
seconds. For example, the Local Mean Time at 9:15 UTC for a launch site whose
longitude is k = 90°45′ East results from
UTC ¼ 9h 15m ¼ ð9 þ 15=60Þ h ¼ 9:25 h

k ¼ 90
450 East ¼ þ ð90 þ 45=60Þ=15 h ¼ þ 6:05 h
LMT ¼ ð9:25 þ 6:05Þ h ¼ 15h 18m
Since the LMT of sunrise and sunset is available or can be computed from
ephemeris data, then the LMT is a convenient quantity for defining launch windows
which depend on local lighting conditions.
In order for a launch to occur in daylight, the earliest launch time is the Local
Mean Time of sunrise at the launch site (subscript L), that is,
LMT LðMINÞ ¼ LMT sunrise at launch site
Likewise, the latest launch time is determined by imposing the condition of

having at least N hours of daylight at the emergency recovery site (subscript E) in
the event of an impact, for the purpose of locating and recovering the vehicle. This
condition can be expressed as follows
LMT EðMAXÞ ¼ LMT sunset at recovery site N h
Assuming negligible motion of the mean Sun within the time interval from
launch to impact, the same condition, in terms of the Local Mean Time at the
launch site, LMTL, can be expressed as follows
DkE
LMT E ¼ LMT L þ
15
In order to fulfil the recovery lighting condition, the latest time of launch must be
DkE
LMT LðMAXÞ ¼ LMT EðMAXÞ
15
and consequently the launch window is the time interval going from LMTL(MIN) to
LMTL(MAX). After choosing the value of LMTL within the time interval indicated
above, the inertial orientation of the orbit is defined by the following expression of
the right ascension of the ascending node, X, as follows

X ¼ aS þ 15 LMT L 12h DkL
where aS is the right ascension of the mean Sun at the time of launch.
second category depend on the additional condition of launching and inserting into
orbit a vehicle in a plane having a specified inertial orientation.
The condition in which the orbital plane of the target vehicle passes through the
launch site is fulfilled only two times per day. For the reason indicated above, the
geocentric latitude of the launch site is assumed to be equal to or less than the
desired inclination of the orbital plane. In practice, it is impossible to launch a
spacecraft at a precise time, due to the finite burn time [45]. Therefore, a small
amount of delta-v must be allowed to execute a plane-change manoeuvre, after
inserting the spacecraft in its parking orbit. Of course, the width of the launch
window depends on the amount of delta-v reserved for plane changes. The equa-
tions governing this case will be shown in a paragraph of this chapter devoted
specifically to rendezvous missions.

third category depend on the additional condition of an orbital plane containing a
specific velocity vector with respect to an inertial system of reference. In order for a
spacecraft launched from the Earth to reach a given planet along a heliocentric
trajectory, it is necessary to consider the heliocentric velocity vector of the Earth
and the hyperbolic excess velocity vector of the spacecraft. A spacecraft arrives at
(or departs from) the sphere of influence of a planet with the so-called hyperbolic
excess velocity vector v∞ (see Sect. 1.7), which is the vector sum of the incoming
(or outgoing) velocity vector of the spacecraft and the velocity of the planet. In the
system of reference whose origin is in the centre of mass of the planet, the direction
(but not the magnitude) of the hyperbolic excess velocity vector changes. The
magnitude, v∞, of the hyperbolic excess velocity vector results from the energy
integral

2 1
v2 ¼ l
r a
when 2/r tends to zero (or the radius vector, r, tends to infinity). This yields
l 12
v1 ¼
a
where a is the negative transverse semi-axis of the hyperbolic trajectory, and l is

the gravitational parameter of the planet.
When a spacecraft goes out of the sphere of influence of the Earth, the hyper-
bolic excess velocity vector, v∞, of the spacecraft is defined by its magnitude, v∞,
and direction (right ascension, a∞, and declination, d∞). The spacecraft, in order to
carry out its mission, must be inserted into an orbital plane which contains the
required hyperbolic excess velocity vector. For this purpose, the spacecraft is
inserted into a circular or elliptical parking orbit around the Earth, and then inserted,
at the proper point of the parking orbit, into an escape hyperbola, in order for the
spacecraft to have the required hyperbolic excess velocity vector v∞ at the moment
of leaving the sphere of influence of the Earth.
An orbital plane containing a specific inertial velocity vector can be obtained
from a given launch site at any time of day, by selecting properly the azimuth angle
of launch. Therefore, the additional requirements concerning lunar or planetary
missions depend on a launch azimuth interval limited only by range safety reasons,
which will be shown in Sect. 4.4. At the present time, let us suppose that such
reasons have determined a minimum allowable value, AL(MIN), and a maximum
allowable value, AL(MAX), for the azimuth angle interval.
In order to define the earliest allowable time of launch, let us consider the
following figure, where the declination, d∞, of the hyperbolic excess velocity
vector is considered to be less than the geocentric latitude, u*L, of the launch site.
Following again Robertson [45], let aL be the right ascension of the launch site,
L. Let aLI be the right ascension of the launch site for the earliest allowable time of
launch. With reference to the preceding figure, a formula of spherical trigonometry
yields
cos i ¼ cos u L sin ALðMINÞ
and the angle Dk is

tan d1
Dk ¼ arcsin
tan i
where 90° < Dk < 180° for the earliest allowable time of launch.
Likewise, the angle DkL is
tan u L
sin DkL ¼
tan i
cos ALðMINÞ
cos DkL ¼
sini
sin DkL
DkL ¼ arctan
cos DkL
Hence, the right ascension of the launch site for the earliest allowable time of
launch is given by
aLI ¼ a1 ðDk DkL Þ
Let k and aS be, respectively, the longitude of the launch site and the right
ascension of the fictitious mean Sun on a particular day. The Universal
Co-ordinated Time of launch corresponding to the computed value of aLI is
UTC LI ¼ aLI k aS þ 12h
As the Earth rotates during the time interval of the launch window, the proper
orientation of the orbital plane is achieved by increasing the azimuth angle of
launch.
The maximum allowable value, AL(MAX), for the azimuth angle interval deter-
mines the latest allowable value for the time of launch, as will be shown below.
By following the same line of reasoning as that indicated above, there results
cos i ¼ cos u L sin ALðMAXÞ

tan d1
Dk ¼ arcsin
tan i
tan u L
sin DkL ¼
tan i
cos ALðMAXÞ
cos DkL ¼
sin i

sin DkL
DkL ¼ arctan
cos DkL
Hence, the right ascension of the launch site for the latest allowable time of
launch is given by
aLF ¼ a1 ðDk DkL Þ
The Universal Co-ordinated Time of launch corresponding to the computed

value of aLF is
UTC LF ¼ aLF k aS þ 12h
The launch window is the period of time, during a given day, which goes from
UTCLI to UTCLF.
Now, let us consider the case in which the declination, d∞, of the hyperbolic
excess velocity vector is greater than the geocentric latitude, u*L, of the launch site.
In this case, the right ascension of the launch site for the earliest allowable time of
launch (aLI) and the related Universal Co-ordinated Time (UTCLI) are given by the
same expressions as those shown above for the case d∞ < u*L.
However, instead of increasing the value of the launch azimuth angle to the
maximum as the Earth rotates, the value of the launch azimuth angle is increased in
this case to less than 90°, according to the geometrical constraint

cos d1
AL ¼ arcsin
cos uL
Then, the value of the launch azimuth angle, AL, decreases with Earth rotation
back to AL(MIN) to define the limit of the launch window.
The same fundamental equations, with suitable quadrant adjustments, are used to
compute the launch window. Because of the geometrical symmetry of this problem,
there are two launch windows per day, which can be determined by means of the
expressions given above. A typical family of launch time curves as a function of
launch azimuth angle is given in the following figure (courtesy of NASA [45]), for
values of d∞ less than, equal to, and greater than u*L.
Of course, the launch window defined by constraints of launch azimuth angle
can be shortened by constraints imposed by lighting conditions or other causes. The
final determination of the launch window must take all such constraints into
account.
4.4 Range Safety
As has been shown in Sect. 4.3, the lowest value of orbital inclination which can be
obtained by launching directly a spacecraft from a given site is equal to the geo-
centric latitude of that site. That minimum value is reached for launches due east
(90° azimuth) or due west (270° azimuth). A satellite launched with an azimuth
between 0° and 180° will have an orbital inclination between 0° and 180°, that is, a
prograde orbit. A satellite launched with an azimuth between 180° and 360° will
have an orbital inclination between 90° and 180°, that is, a retrograde orbit. In order
to obtain a polar orbit, it is necessary to launch due north (0° azimuth) or due south
(180° azimuth), independently of the geographic position of the launch site. Since
the Earth spins to east at a mean linear velocity of 465.1 m/s at the equator [48] and
0 m/s at the poles, substantial savings in rocket propellant can be obtained by
locating a launch site on or near the equator and launching due east. Conversely, in
order to insert a satellite into a retrograde orbit, it is desirable to launch it from a site
located at as high a latitude as possible, because the launcher has to overcome as
little of the Earth rotational velocity as possible. Such constraints as those indicated
above are only due to geometric and energetic reasons.
There are additional constraints, which are due to safety reasons. In particular, it
is necessary to consider the possibility of an impact between the vehicle (or the
booster stage) and the ground. This impact may occur in a populated or high-traffic
region. This problem can be solved by imposing limits to the azimuth angle
interval, when the launch occurs from a given site.
When the boost trajectory is within such limits, there is a very low probability of
stage or vehicle impact in populated areas. In the event of the vehicle drifting
outside the safety limits, then the vehicle is usually destroyed.
The following figure, due to the courtesy of the National Security Space Institute
[49], shows the security limits for launch azimuth and the consequent orbital
inclination angles, in case of launches from Cape Canaveral Air Force Station and
Kennedy Space Centre in Florida. As shown in the figure, security limits for launch
azimuth angles cause limits in orbital inclinations, i, which can be achieved when a
launch is performed from a given site.
As has been shown in Sect. 4.3, the formula of spherical trigonometry, which
expresses the orbital inclination, i, as a function of the geocentric latitude, u*L, of
the launch site and the launch azimuth, AL, is
4.4 Range Safety 621
cos i ¼ cos u L sin AL
where u*L i < 90° for prograde orbits. The preceding equation, solved for AL,
is

cos i
AL ¼ arcsin
cos u L
which shows mathematically the reason for which i cannot be less than u*L.
Otherwise, if i < u*L, then the argument, (cos i)/(cos u*L), of the inverse sine
function would be greater than unity, which is impossible. Therefore, the value
AL = 90° (launch due east) of the launch azimuth corresponds to the orbit of
minimum inclination (iMIN = u*L) which can directly be achieved by launching a
spacecraft from a site of geocentric latitude u*L. In the case of Cape Canaveral
(u*L 28°.5N), the preceding figure shows that the azimuth angle interval
imposed by safety reasons is between 35° and 120°. Therefore, by launching a
spacecraft from Cape Canaveral at the azimuth angle nearest 90° (which, in the
present case, is just 90°), it is possible to achieve an orbit of minimum inclination
(iMIN = 28°.5), and by launching at the azimuth angle furthest from 90° (which, in
the present case, is 35°), it is possible to achieve an orbit of maximum inclination
(iMAX = 59°.7). The azimuth angle, AL, indicated above and the consequent orbital
inclination, i, are given in an inertial system of reference. However, the azimuth
angle measured by means of a compass in the point of the surface of the Earth
where the launch takes place is given by the instrument in a system of reference
which rotates with the Earth. It is necessary to compensate for this rotation. To this
end, it is also necessary to compute the value of the desired azimuth angle, ALrot,
with respect to a rotating system of reference.
For example, let us suppose that we want to insert an artificial satellite into a
circular orbit of altitude h = 350 km. The orbital radius corresponding to this
altitude is r = (6378.137 + 350.000) km = 6728.137 km. The magnitude of the
orbital velocity vector results from the energy integral
1 1
2 1 2 lE 12 398600:44 2
lE ¼ ¼ ¼ 7:697 km/s
r a r 6728:137
Let us consider the case of a prograde orbit (i < 90°), a desired orbital inclination
i = 51°.6494, which is the orbital inclination of the International Space Station on
14 October 2013, at 10:18 UTC [50], and a launch from Cape Canaveral (geo-
centric latitude u*L 28°.5N). The inertial azimuth angle results from the formula
given above

cos i cos 51
:6494
AL ¼ arcsin ¼ arcsin ¼ 44
:9129
cos u L cos 28
:5
At the equator, the Earth rotates at a linear velocity (vE)eq = 0.4651 km/s [48].
Therefore, the linear velocity of the Earth at the site of launch is
vE ¼ ðvE Þeq cos u L ¼ 0:4651 cos 28

:5 ¼ 0:4087 km=s
With reference to the following figure, let v and vE be the velocity vectors of,
respectively, the satellite and the Earth with respect to an inertial system.
Let AL be the azimuth of launch with respect to an inertial system of reference.

The value of this angle, in the example given above, has been found to be 44°.9129.
Let vrot and ALrot be, respectively, the velocity vector of the satellite and the azimuth
of launch with respect to a system of reference rotating with the Earth. The angle
ALrot is the azimuth measured by an observer located in the launch site. As shown in
the following figure, there results v = vE + vrot.
It is necessary to determine the magnitude, vrot, and the azimuth angle, ALrot, of
the velocity vector, vrot, of the satellite with respect to a system of reference rotating
with the Earth and having its origin in the launch site. To this end, we determine the
horizontal component, (vrot)x, and the vertical component, (vrot)y, of the vector vrot.
By inspection of the preceding figure, these components result from
ðvrot Þx ¼ v sin AL ðvE Þeq cos u L ¼ ð7:697 sin 44

:9129 0:4651 cos 28
:5Þ km/s
¼ 5:02558 km=s
ðvrot Þy ¼ v cos AL ¼ ð7:697 cos 44
:9129Þ km=s ¼ 5:45087 km=s
Hence, the azimuth angle, ALrot, of the velocity vector, vrot, of the satellite with
respect to a system of reference rotating with the Earth results from
" #
ðvrot Þx 5:02558
ALrot ¼ arctan ¼ arctan ¼ 42
:6754
ðvrot Þy 5:45087
4.4 Range Safety 623
and the magnitude, vrot, of the same velocity vector results from
h i12 1
vrot ¼ ðvrot Þ2x þ ðvrot Þ2y ¼ 5:025582 þ 5:450872 2 ¼ 7:41407 km/s
The preceding value refers to the velocity which the rocket must impart to the
satellite in order to insert it into the desired orbit. Since the magnitude of the orbital
velocity vector v is v = 7.697 km/s, then the difference
v vrot ¼ ð7:697 7:414Þ km=s ¼ 0:283 km=s
is the velocity given to the satellite at the moment of launch by the rotating Earth.
A launch performed in the manner shown above inserts a satellite into an orbit of
altitude h = 350 km on the surface of the Earth and inclination i = 51°.65 with
respect to the equatorial plane.
When the desired orbital inclination is not compatible with the conditions
imposed by either the geocentric latitude of the launch site or the permitted azimuth
interval, then it is necessary to:
• insert the satellite into an initial orbit whose inclination satisfies such conditions;
and
• perform successively a manoeuvre of plane change.
The plane-change manoeuvre must be performed at one of the two nodes of
intersection between the initial orbit and the final orbit, unless it is desired to change
not only the inclination angle but also other orbit elements of the initial orbit. In the
latter case, a two-impulse manoeuvre is needed.
There are several launch sites around the world. Those of them which are most
frequently used for inserting satellites into orbits are shown in the following table
(from Ref. [51]), which indicates the geodetic latitude, longitude, minimum and
maximum allowed azimuth for each site.
Country Launch site uL kL ALrot(MIN) ALrot(MAX)

USA Cape Canaveral 28°.5N 81°.0W 35° 120°
USA Vandenberg 34°.4N 120°.35W 147° 201°
Russia/CIS Baikonur 45°.6N 63°.4E 65° −13°
Russia/CIS Plesetsk 62°.8N 40.1E – –
France/Europe Kourou 5°.2N 52°.8W −11° 90°
Japan Tanegashima 30°.4N 130°.9E 0° 180°
China Jiuquan 40°.6N 99°.9E – –
China Xichang 28°.3N 102°.0E – –
India Sriharikota 13°.8N 80°.3E – 140°
4.5 Ascent Trajectories
According to Robertson [45], there are fundamentally two types of boost trajec-
tories for orbital missions. They are:
(a) powered ascent plus coasting; and
(b) direct ascent with continuous burn.
A trajectory of the first type is similar to that of a short-range or medium-range
ballistic missile. The vehicle goes firstly through the phase of main powered ascent,
and then through the phase of coasting along an elliptical non-powered path, at
whose apogee it receives a short impulse in order to acquire the necessary orbital
velocity. Further impulses may be given for orbit raising, as will be shown in the
following paragraphs.
A trajectory of the second type results from a continuous burn, except during
staging, with vehicle guidance from the moment of launch to the moment of
insertion into orbit, at which time the thrust ceases. As a general rule, the vehicle is
launched vertically and then pitched after a few seconds in the direction of the
required azimuth. The pitching is performed by tilting the thruster, according to one
of the methods described in Sect. 4.2.
The ascent trajectory is determined by integrating numerically the equations of
motion, in the time interval going from launch to orbit insertion, as will be shown
below. These equations result from equating the time derivative of momentum and
the sum of all the forces Fi (i = 1, 2, … , N) acting on the vehicle, that is,
dðmvÞ X N
¼ Fi
dt i¼1
The principal forces which act on the vehicle are thrust, aerodynamic forces,
gravitational force, centrifugal force, and Coriolis force. The latter two of these are
fictitious forces, which are taken into account because the motion of the vehicle is
considered with respect to a non-inertial system of reference which rotates with the
Earth.
The first of the forces considered here is the thrust provided by the vehicle. As
has been shown in Sect. 4.2, the magnitude of the thrust is
F ¼ mv
where m_ is the mass flow rate of propellant, ve is the gas velocity at the exit plane of
the nozzle, Ae is the area of the cross section at the same plane, pe is the gas pressure
at the same plane, and p0 is the ambient pressure.
The total impulse, IT, results from integrating the thrust over the total burning
time, tT, as follows
4.5 Ascent Trajectories 625
ZtT
IT ¼ Fdt
0
Then, the specific impulse, Isp, which is the impulse provided per unit weight of
propellant or the thrust per unit weight of propellant burned in the unit time, results
from
IT F
Isp ¼ ¼
WT W_
where WT is the total weight of propellant burned, and W _ is the propellant weight
flow rate.
The specific impulse is used to compute the ideal velocity gain, DvI, which may
be obtained (neglecting all the sources of loss) from an impulse. The ideal velocity
gain is given by Tsiolkovsky’s equation (see Sect. 3.25):

mprop þ mpay
DvI ¼ Isp g0 ln
mpay
where g0 is the acceleration of gravity at sea level, mprop is the propellant mass, and
mpay is the payload mass (that is, the residual mass after burn).
Aerodynamic forces are very important in the early phase of boost. As has been
shown in Sect. 3.19, such forces depend on the vehicle shape, the atmospheric
density, the air velocity relative to the vehicle, and the projected area of the vehicle
normal to the relative velocity vector. The aerodynamic forces are sometimes
expressed in terms of the dynamic pressure, q, which is defined as follows
1 2
q¼ qv
2 rel
where q is the atmospheric mass density, and vrel is the air velocity relative to the
vehicle. As is well known, the component of the aerodynamic force parallel to the
relative velocity vector, vrel, is called the drag force, D, and the component per-
pendicular to the relative velocity vector is called the lift force, L, which are
expressed by
D ¼ qACD
L ¼ qACL ¼ qACL0 a
These forces, if expressed with respect to the vehicle axes, are the axial force, X,
and the normal force, N, which are given by
X ¼ qACX
N ¼ qACN ¼ qACN0 a
where A is the projected area of the vehicle normal to the relative velocity vector,
CD is the drag coefficient, CL is the lift coefficient, a is the angle of attack, C′L is the
derivative of the lift coefficient with respect to the angle of attack, CX is the axial
force coefficient, CN is the normal force coefficient, and C′N is the derivative of the
normal force coefficient with respect to the angle of attack.
The weight, W, of the vehicle is the gravitational force exerted by the Earth on
the vehicle. The weight of the vehicle varies with time, because of the propellant
consumption during the boost phase. Therefore, W decreases according to
_ ðt t0 Þ
W ¼ W0 W
Likewise, the mass, m, of the vehicle decreases according to
m ¼ m0 m_ ðt t0 Þ
where W0 and m0 are, respectively, the initial weight and the initial mass of the
vehicle, t and t0 are, respectively, the current time and the initial time, and W_ and m_
are the flow rates of, respectively, the weight and the mass of the propellant.
The weight of the vehicle depends not only on the propellant consumption, but
also on the vehicle altitude over the Earth, according to Newton’s law of gravitation
l
W ¼ mg ¼ m
r2
where l is the gravitational parameter of the Earth, and r is the radius vector of the
vehicle.
The Coriolis force and the centrifugal forces must be taken into account because
the motion of the vehicle is considered with respect to a non-inertial system of
reference which rotates with the Earth. The Coriolis force vector, C, and the cen-
trifugal force vector, CF, are expressed as follows
C ¼ 2m xE vrot
CF ¼ m xE ðxE rrot Þ
where xE is the angular velocity vector of the Earth with respect to an inertial
geocentric system of reference, rrot and vrot are, respectively, the position and the
velocity vectors of the vehicle in the rotating system of reference, and m is the mass
of the vehicle. The principal forces acting on the vehicle are illustrated in the
following figure. We consider here only the motion of the centre of mass of the
vehicle. The moments M and their effects on the motion of the vehicle will be
considered in Chap. 9.
Other forces smaller than those indicated above, such as thrust misalignment,
unsymmetrical aerodynamic shape of the vehicle, non-spherical shape of the Earth,
third-body gravitational forces, etc., are not considered here.
The second principle of dynamics is expressed by the following equation
m_vrot ¼ F þ D þ L þ W 2m xE vrot mxE ðxE rrot Þ
In order to solve this equation, it is convenient to use a system of reference

whose co-ordinates are the tangential direction and the normal direction of flight.
The tangential direction of flight is the same as that of the velocity vector, vrot, with
respect to the rotating system of reference, and the normal direction of flight is
perpendicular to it. During the ascent through the low atmospheric layers, the Earth
is assumed to be stationary. The rotation of the Earth is taken into account after the
vehicle has reached a height where the effect of the aerodynamic forces on the
trajectory can be neglected. The correction due to the rotation of the Earth in the
early stage of ascent will subsequently be taken into account.
The preceding equation, written in terms of tangential and normal co-ordinates,
gives rise to the following scalar equations
m_vrot ¼ F cos arot D W cos hrot
mvrot h_ rot ¼ F sin arot þ L þ W sin hrot mvrot w_ rot
As shown in the preceding figure, the angle hrot is the conjugate flight path angle
between the local vertical and the velocity vector, vrot; the angle wrot is the angle
between the position vector of the vehicle at the moment of launch and the position
vector at the current time (t). The term containing the time derivative of wrot is the
fictitious centrifugal force. The differential of wrot results from
vrot dt sin hrot

dwrot ¼
rrot
hence, the time derivative of wrot is
dw vrot sin hrot

w_ rot rot ¼
dt rrot
The preceding expression, introduced into the equations of motion written

above, yields
F cos arot D
v_ rot ¼ g cos hrot
m m

F sin arot L g vrot
h_ rot ¼ þ þ sin hrot
mvrot mvrot vrot rrot
The velocity and the conjugate flight path angle of the vehicle result from a first
integration of these equations, as follows
Zt
vrot ¼ v_ rot ds
0
Zt
hrot ¼ h_ rot ds
0
The altitude and ground range result from a further integration, as follows
Zt
hrot ¼ vrot cos hrot ds
0
Zt
vrot sin hrot
srot ¼ r0 ds
rrot
0
In the general case, it is necessary to use numerical methods to perform these

integrations.
As Robertson [45] points out, the differential equation of velocity can be inte-
grated analytically by neglecting the drag term and making other simplifying
assumptions. The result is useful for the purpose of expressing the velocity as a
function of the ideal velocity gain, DvI, and the gravity loss. This makes it possible
to evaluate the thrusting manoeuvres which are necessary near the Earth. To this
end, let us assume a constant or mean value for the thrust, F, gravitational accel-
eration, g, and conjugate flight path angle, hrot. The angle of attack, arot, of the
vehicle is assumed to be zero, so that F cos arot = F. In such conditions, the
trajectory is rectilinear, and the differential equation of velocity can be integrated as
follows
Zt Zt
F
vrot ¼ ds g cos hrot ds
m
t0 t0
By substituting
F ¼ Isp W_
1
m¼ W0 W_ ðt t0 Þ
g0
mprop þ mpay W0
uv ¼
mpay _ ðt t0 Þ
W0 W
into the expression of vrot and integrating, there results
vrot ¼ ðvrot Þ0 þ Isp g0 lnðuv Þ ðg cos hrot Þðt t0 Þ
where F is the assumed average value of the thrust, Isp is the assumed average value
of the specific impulse, m is the mass of the vehicle at the current time (t), g0 is the
acceleration of gravity at sea level, W0 is the initial weight of the vehicle, W_ is the
weight flow rate of the propellant, t and t0 are, respectively, the current time and
the initial time, g is the assumed average value of the gravitational acceleration, and
(vrot)0 is the initial velocity of the vehicle in the rotating system of reference.
The second term and the third term on the right-hand side of the preceding
equation are, respectively, the ideal velocity gain, DvI, and the gravity loss.
The ideal velocity gain, when a multi-stage vehicle is used, can be computed by
adding the ideal velocity gains of the single stages, with adjustments to the pro-
pellant mass ratio, as follows
DvI ¼ Isp1 g0 lnðuv1 Þ þ Isp2 g0 lnðuv2 Þ þ
The computation shown above has been performed with respect to a system of
reference rotating with the Earth. Consequently, the velocity, vrot, of the vehicle,
whose expression has been given above, is the velocity in a reference system
rotating with the Earth, not the velocity in an inertial system.
As has been shown in Sect. 4.4, in order to compute the velocity vector, v, of the
vehicle with respect to an inertial system, it is necessary to add the linear velocity
vector, vE, of the Earth at the site of launch to the velocity, vrot, of the vehicle with
respect to a system of reference rotating with the Earth, as follows
v ¼ vE þ vrot
Let v and h be, respectively, the magnitude of the velocity vector v and the
conjugate flight angle path with respect to an inertial system of reference.
By inspection of the preceding figure, there results

h i2 12
v¼ ðvE Þeq cos uL þ vrot sin ALrot þ v2rot cos2 ALrot
h i12
¼ ðvE Þ2eq cos2 uL þ 2ðvE Þeq vrot cos uL sin ALrot þ v2rot
v
rot
h ¼ arccos cos hrot
v
where (vE)eq = 0.4651 km/s [48] is the linear velocity of the Earth at the equator,
u*L is the geocentric latitude of the launch site, vrot is the velocity of the vehicle in
the rotating reference system, ALrot is the launch azimuth measured in the rotating
reference system, and hrot is the conjugate flight path angle in the rotating reference
system.
4.6 Insertion into Orbit
The inertial velocity, vins, which a vehicle must possess at the moment of its
insertion into orbit is important for the purpose of estimating the performance
requirements of the launcher. The following expression (see Sect. 4.5)
4.6 Insertion into Orbit 631
DvI ¼ Isp1 g0 lnðuv1 Þ þ Isp2 g0 lnðuv2 Þ þ
of the ideal velocity gain for a multi-stage launcher can be used for a rough
estimation.
Let a spacecraft be inserted into an elliptic orbit of radii rP = (RE)eq + hP at
perigee and rA = (RE)eq + hA at apogee, where (RE)eq is the equatorial radius of the
Earth. Let the perigee be the point of insertion. The vis-viva integral is

2 1
v2 ¼ lE
r a
where lE and a are, respectively, the gravitational parameter of the Earth and the
major semi-axis of the elliptic orbit. By substituting r = rP and a = ½(rP + rA) and
solving for v, there results
12
2lE rA
vins ¼
rP ðrP þ rA Þ
In the particular case of a circular orbit of height h, then

rP = rA = r = (RE)eq + h, and the preceding equation reduces to
" #12
l 12 lE
vins ¼ E
¼
r ðRE Þeq þ h
The cases shown above concern the insertion of a spacecraft into orbit at perigee.
Let us consider now the general case in which the initial conditions of a rocket at
the point P0 of burnout are r = r0, v = v0, and c = c0, where c is the flight path angle
(c = p/2 − h) measured outward from the local horizontal to the local velocity
vector, v0, as shown in the following figure. We want to determine the eccentricity, e,
of the orbit and the true anomaly, /0, of the point P0 of burnout.
Following Thomson [52], let rP and r0 and be the radii of, respectively, perigee
(/ = 0) and burnout point (/ = /0).
At perigee, the equation of a conic section in polar co-ordinates
p h2 1
r¼ ¼
1 þ e cos / lE 1 þ e cos /
yields
h2
¼ r P ð 1 þ eÞ
lE
Hence, the equation of a conic section may be written as follows
r P ð 1 þ eÞ
r¼
1 þ e cos /
By inspection of the preceding figure, the transverse (or horizontal) and radial
(or vertical) components of the initial velocity vector, v0, are

d/ h
v0 cos c0 ¼ r0 r0 /00 ¼
dt 0 r0

dr 0 dr d/ l e sin /0
v0 sin c0 ¼ r ¼ ¼ E
dt 0 0 d/ 0 dt 0 r0 v0 cos c0
The preceding equation, solved for e sin /0, yields
r0 v20
e sin /0 ¼ sin c0 cos c0
lE
The general equation of a conic section, written for r = r0, is
p h2 1
r0 ¼ ¼
1 þ e cos /0 lE 1 þ e cos /0
With reference to the preceding figure, the magnitude of the moment of

momentum per unit mass of the spacecraft is by definition
h ¼ rv sin h ¼ rv cos c
Therefore
h2 ¼ r 2 v2 cos2 c
By inserting the preceding expression into the equation of r0, there results
r02 v20 cos2 c0 1

r0 ¼
lE 1 þ e cos /0
By dividing the terms on both sides of the preceding equation by r0, there results
r0 v20 cos2 c0
¼ 1 þ e cos /0
lE
The preceding equation, solved for e cos /0, yields
r0 v20
e cos /0 ¼ cos2 c0 1
lE
By dividing the expression of e sin /0 by the expression of e cos /0, there results
2
r0 v0 =lE sin c0 cos c0
tan /0 ¼ 2
r0 v0 =lE cos2 c0 1
By adding the square of e sin /0 to the square of e cos /0, there results
2 2
r 0 v0
e ¼
2
1 cos2 c0 þ sin2 c0
lE
The last two of the equations written above define completely the orbit (in terms
of /0 and e) of the spacecraft for any initial conditions expressed by r0, v0, and c0 at
burnout. The mechanical energy per unit mass of the spacecraft results from
summing its potential energy (−lE/r) to its kinetic energy (½v2), as follows
lE 1 2
E¼ þ v
r 2
The preceding equation, written in non-dimensional form for the point of

burnout (r = r0 and v = v0), is
Er0 1 r0 v20
¼ 1
lE 2 lE
As a result of the equations written above, in case of c0 different from zero, the
orbital eccentricity, e, can never be zero, that is, the orbit can never be circular.
In addition, when (r0v20/lE)cos2c0 = 1, then /0 = 90°. When (r0v20/lE)cos2c0 < 1
and c0 > 0, then /0 is in the second quadrant.
When a high-altitude orbit is desired, the spacecraft is usually inserted into an
initial low-altitude circular orbit, followed by an elliptic transfer orbit to the desired
orbital radius, followed in turn by a circularisation of the elliptic orbit. For this
purpose, a Hohmann transfer orbit, illustrated in the following figure, is frequently
used. As will be shown at length in the next paragraphs, a Hohmann transfer orbit
(so called after the German engineer Walter Hohmann [53]) is the two-impulse
trajectory which uses the least amount of propulsive energy.
A Hohmann transfer orbit between two circular orbits of radii, respectively,

r1 and r2 (such that r1 < r2) is an elliptic orbit whose perigee radius is rP = r1
and whose apogee radius is rA = r2. By imposing this condition, it is possible to
determine the major semi-axis, a, and the eccentricity, e, of a Hohmann transfer
orbit as follows
r1 þ r2
a¼
2
r1 r2 r1
e¼1 ¼
a r2 þ r1
When the spacecraft is inserted into a Hohmann transfer orbit at perigee, as

shown in the preceding figure, then the insertion velocity results from the expres-
sion given above
12
2lE rA
vins ¼
rP ðrP þ rA Þ
By contrast, when the spacecraft is inserted into a Hohmann transfer orbit at any
other point (of radius rins > rP) than perigee, then the insertion velocity depends on
the insertion radius, rins, insertion flight path angle, cins, with respect to the local
horizontal, and apogee radius, rA, as indicated by Robertson [45]:
" #12
2lE ðrA rins ÞrA
vins ¼
rins rA2 rins
2 cos2 c
ins
The spacecraft velocity at the apogee of a Hohmann transfer orbit is [45]:

rins
vA ¼ vins cos cins
rA
Therefore, the total velocity, vT, required to insert the spacecraft into its final
circular orbit of radius r2 = rA is [45]:
12
lE rins
vT ¼ vins þ vins cos cins
rA rA
4.7 Rendezvous Manoeuvres
A rendezvous manoeuvre involves two space vehicles which must not only reach
the same orbit but also approach at a very close distance. One of such vehicles,
called the target, is passive and does not perform the manoeuvre. An example of
target is the International Space Station, which revolves about the Earth. The other
vehicle, called the interceptor or the ferry or the shuttle, is active and performs the
manoeuvre necessary to place itself within visual contact with the target. An
example of an interceptor is the Space Shuttle or another space vehicle designed to
reach the Space Station. For this purpose, the position and velocity vectors of the
interceptor must be made to match with those of the target. After completing a
rendezvous manoeuvre, the distance between the two vehicles will be kept constant
by performing station-keeping manoeuvres.
A rendezvous manoeuvre is sometimes followed by docking or berthing
manoeuvres, during which the two vehicles come into physical contact and
establish a link between them.
A rendezvous manoeuvre is more complex than one for a transfer trajectory from
one orbit to another, because it must be performed under severe time constraints.
There are two categories of rendezvous manoeuvres. The first category includes
the manoeuvres performed by the interceptor coming directly from the ground to
reach the orbiting target. For this purpose, the ascent of the interceptor may start
either directly from a point of the ground or from a point of an intermediate orbit.
The second category includes the manoeuvres performed by the interceptor
revolving about the Earth in an arbitrary orbit to place itself on the same orbit as
that of the target. A particular case is that in which the interceptor and the target
occupy two different positions of the same orbit.
Let us consider now the rendezvous manoeuvres belonging to the first category.
In case of a direct ascent of the interceptor from a launch site placed on the surface
of the Earth, the ascent trajectory must lie on the same plane as that of the orbit of
the target. This fact imposes severe constraints on the time of launch, because the
launch site is located in a point of the rotating Earth, and this point passes through
the orbital plane of the target only two times per day. In addition, at each of these
two favourable times, the target must be in such a position along its orbit as to be
reached by the interceptor.
There are two manners to attain this object. The first manner consists in inserting
the target not into an arbitrary orbit but into a rendezvous-compatible orbit,
meaning by this name an orbit in which the target passes over the launch site at least
once per day. A target inserted into such an orbit can be made to pass over the
launch site once or twice per day. By so doing, there are several opportunities for
direct ascent rendezvous manoeuvres over the total duration of a mission.
The constraints posed on the launch time, tL, in case of a direct ascent of the
interceptor without plane change are illustrated in the following figure.
4.7 Rendezvous Manoeuvres 637
The Greenwich Hour Angle, GHAL, of the launch site, L, required in order for
the launch site to be in the same plane as that of the target vehicle results from
GHAL ¼ Xt þ DkL kL
where

tan u L
DkL ¼ arcsin
tan it
By the way, the Greenwich Hour Angle of a celestial body is the angular
distance of that body, measured westward along the celestial equator, from the
Greenwich celestial meridian.
In case of a launch site located in the northern hemisphere, the difference in
longitude, DkL, is such that
0
DkL 90
for northerly launches

90 \DkL \180 for southerly launches
u*L and kL are, respectively, the geocentric latitude and the longitude of the launch
site, it is the inclination with respect to the equator of the orbital plane of the target
vehicle, Xt = Xt0 + X′t TD is the current value (including the correction due to the
oblateness of the Earth) of the right ascension of the ascending node of the target
orbit, Xt0 is the right ascension of the ascending node of the target orbit at the
moment, t0, of the insertion of the target vehicle into its orbit, and TD is the time
interval going from the time t0 to the current time t. In addition, as has been shown
in Sect. 3.2, the regression (for a prograde orbit) rate of the right ascension of the
ascending node is
2
3 rE cos it
X0t ¼ nt C20 2
2 at 1 e2 t
where rE is the equatorial radius of the Earth, at, et, and nt are, respectively, the
major semi-axis, the eccentricity, and the mean motion of the orbit of the target, and
finally C20 = −0.001082636 is the principal zonal coefficient of the gravitational
field of the Earth.
Let GHA0 be the Greenwich Hour Angle at some epoch, which may for con-
venience be taken at the time of the initial perigee passage of the target. According
to Robertson [45], the times of launch of the interceptor with respect to the epoch
are given by
GHAL GHA0 þ 2m p
tL ¼
xE
where m is the whole number of days elapsed since epoch, and xE is the angular
velocity at which the Earth rotates about its axis.
The equation given above defines the times, measured with respect to the epoch,
at which the interceptor may be launched to reach the orbital plane of the target.
The same equation does not define the position of the target with respect to the
point of launch of the interceptor. Therefore, a further equation is needed to
determine whether a rendezvous between the two vehicles can occur for a launch
opportunity.
To this end, the times of launch of the interceptor with respect to the epoch may
also be expressed as follows
tL ¼ DtR tascent tboost þ nTt
where DtR is the time interval going from the perigee passage to the rendezvous
point along the orbit of the target, tascent is the time interval going from the burnout
to the rendezvous point along the trajectory of the interceptor, tboost is the burning
time in the boost trajectory of the interceptor, n is the whole number of revolutions
of the target from the epoch to perigee passage immediately before rendezvous, and
Tt = 2p(a3t /lE)½ is the orbital period of the target.
Of course, the times computed by means of the last two equations written above
must be equal for a direct ascent of the interceptor leading to a successful
rendezvous.
The launch time computed by means of tL = (GHAL − GHA0 + 2 mp)/xE yields
a fixed set of values for any given launch site and target orbit. On the other hand,
equation tL = DtR − tascent − tboost + nTt contains quantities which are independent
of the launch site and target orbit. Since such quantities are variable within some
limits, then their values can be chosen so as to lead to values of tL compatible with
those resulting from tL = (GHAL − GHA0 + 2 mp)/xE.
In particular, the quantity tboost, which depends on the booster vehicle and on the
boost guidance techniques, is scarcely variable. The same holds with the quantity
nTt, which depends on the number of revolutions of the target within the time
interval in which a rendezvous is attempted. Therefore, the available quantities are
tascent and DtR. The latter, which is the time interval going from the perigee passage
to the rendezvous point along the target orbit, may be held to depend on the time
interval tascent going from the interceptor burnout to the rendezvous point, as the
sequel will show. Since the time of launch for a given opportunity results from
tL = (GHAL − GHA0 + 2 mp)/xE, and since it is advantageous to insert the inter-
ceptor into its ascent trajectory near perigee, then the orientation of the line of
apsides of the ascent trajectory can be considered defined. Therefore, there is a
family of ascent ellipses having all of them the same radius of perigee (rP) and the
same argument of perigee (x), but different radii of apogee (rA), depending on the
value of the velocity increment (Dv) given by the launcher to the interceptor to
insert it into its ascent trajectory. To each value of velocity increment, there cor-
responds a particular point in which the ascent trajectory of the interceptor
intersects the orbit of the target. Therefore, the time of intercept, DtR, can be
considered as a function of the time of ascent, tascent. The latter, in turn, can be used
as an independent variable in an iterative process aimed at determining values of tL
such as to satisfy both of the following equations
GHAL GHA0 þ 2m p
tL ¼
xE
This iterative process is explained below. Let us suppose to know the value of
the velocity, vins, at which the interceptor is inserted into its ascent trajectory.
The major semi-axis, as, of the ascent trajectory results from the vis-viva integral
v2/l = (2/r − 1/a) as follows
lE ðRE þ hins Þ
as ¼
2lE ðRE þ hins Þv2ins
where hins and vins are, respectively, the insertion altitude and velocity of the
interceptor, and RE and lE are, respectively, the equatorial radius and the gravita-
tional parameter of the Earth.
The eccentricity, es, of the ascent trajectory results from 1 − e2 = p/a = h2/(la)
(see Sect. 1.1). The preceding equation, solved for e, yields
12
h2
es ¼ 1 s
l E as
where h2s = hs hs = (rs vs) (rs vs) is the squared magnitude of the moment
of momentum per unit mass in the motion of the interceptor. As has been shown in
Sect. 1.1, h is a constant of motion in the two-body problem, in accordance with
Kepler’s second law. In addition, when the interceptor is inserted into its ascent
trajectory at perigee (in which point the flight path angle, c, is equal to zero), then
h2s = (rinsvins)2 = [(RE + hins)vins]2, where hins is the insertion altitude of the inter-
ceptor, and RE is the equatorial radius of the Earth. In the general case, there results
h2s = (rinsvins cos cins)2, where cins is the flight path angle.
According to Robertson [45], the argument of perigee, xs, of the ascent tra-
jectory results from

sin u L
xs ¼ arcsin þ hboost
sin it
where, with reference to the following figure and assuming a launch site, L, located
in the northern hemisphere, hboost is the boost angular range, and xs is contained in
the following interval
0
\xs \90

90 \xs \180 for southerly launches
Let xt and xs be the arguments of perigee of, respectively, the target orbit and
the ascent trajectory.
As shown in the preceding figure, the angular displacement, Dx, between the
two apsidal lines is
Dx ¼ /s /t
Let /t and /s be the true anomalies of one of the two interception points (I1 and I2)
measured from, respectively, the perigee of the target orbit and the perigee of the
ascent trajectory. The true anomaly /t results from /t = /s − Dx.
The radius vector, r, of the interception point results from

at 1 e2t as 1 e2s
r¼ ¼
1 þ et cos /t 1 þ es cos /s
By substituting cos /t = cos(/s − Dx) = cos /s cos Dx + sin /s sin Dx into

the expression of r written above, there results

pt pt
et cos D x es cos /s þ ðet sin D xÞ sin /s ¼ 1
ps ps
where pt = at(1 − e2t ) and ps = as(1 − e2s ). The quantities within parentheses in the
preceding equation do not depend on /s. By setting for convenience
pt pt
A ¼ et cos D x es B ¼ et sin D x C¼ 1
ps ps
the preceding equation becomes
A cos /s þ B sin /s ¼ C
By squaring the terms on both sides of the preceding equation, remembering the
trigonometric identity sin2/s = 1 − cos2/s, and taking into account that
2AB sin /s cos /s ¼ ð2A cos /s ÞðB sin /s Þ ¼ ð2A cos /s ÞðC A cos /s Þ
because A cos /s + B sin /s = C, the preceding equation becomes

A2 þ B2 cos2 /s 2AC cos /s þ C 2 B2 ¼ 0
By solving the equation for cos /s, there results

1 1
AC ½A2 C 2 ðA2 þ B2 ÞðC2 B2 Þ2 AC BðA2 þ B2 C2 Þ2
cos /s ¼ ¼
A2 þ B2 A2 þ B2
There are two solutions, depending on which sign (plus or minus) is chosen in
front of B. The solution corresponding to the plus sign is
1
BC AðA2 þ B2 C 2 Þ2
sin /s ¼
A2 þ B2
1
AC þ BðA2 þ B2 C 2 Þ2
cos /s ¼
A2 þ B2
The solution corresponding to the minus sign is

1
BC þ AðA2 þ B2 C2 Þ2
sin /s ¼
A2 þ B2
1
AC BðA2 þ B2 C 2 Þ2
cos /s ¼
A2 þ B2
There are either two distinct points (I1 6¼ I2), or two coincident points (I1 = I2),
or no points of intersection between the two elliptical orbits, depending on whether
the value of A2 + B2 − C2 be, respectively, greater than, or equal to, or less than
zero. Remembering the following expressions of Sect. 1.3
the eccentric anomaly, Æs, corresponding to the true anomaly, /s, computed above
results from
Since the interceptor is inserted at perigee into its ascent trajectory, then tP = 0,
and the ascent time is
12
Ms a3s
tascent ¼ ¼ Ms
ns lE
This value of tascent is to be substituted into tL = DtR − tascent − tboost + nTt.

The true anomaly of the interception point in the target orbit results from the
preceding equation /t = /s − Dx, because /s and Dx are known.
The interval of time, DtR, going from perigee to the interception point in the
target orbit results from the following equations
12
Mt a3t
DtR ¼ ¼ Mt
nt lE
The iterative process shown above can be summarised as follows. The velocity,
vins, at which the interceptor is inserted into its ascent trajectory is assumed as an
independent variable. Supposing to know the value of vins, the equations given
above make it possible to compute sequentially the values of as, es, xs, Dx, /t, pt,
ps, A, B, C, cos /s, sin /s, cos Æs, sin Æs, Æs, Ms, tascent, cos Æt, sin Æt, Æt, Mt, and
DtR. Now, the values of tascent and DtR computed above are substituted into the
following equation
where n is the whole number of revolutions of the target from the epoch to perigee
passage immediately before rendezvous, and Tt = 2p(a3t /lE)½ is the orbital period
of the target. This equation expresses the time of launch of the interceptor with
respect to the epoch.
Now, the value of tL resulting from the preceding equation is compared with the
set of values resulting from
GHAL GHA0 þ 2mp

tL ¼
xE
where m is the whole number of days elapsed since the epoch, and xE is the angular
velocity at which the Earth rotates about its axis. This comparison is aimed at
ascertaining the compatibility between the values of tL resulting from the two
equations indicated above. When these values are found compatible, the process
comes to an end. Otherwise, it is necessary to assign a new value to vins and perform
a further iteration.
As has been shown above, it is not easy to determine a suitable ascent trajectory
for a rendezvous between an interceptor and a target, when the orbit of the latter is
arbitrary. Therefore, when a such a manoeuvre is necessary, it is convenient to
insert the target into a rendezvous-compatible orbit. Given the performance of the
boost vehicle, a rendezvous-compatible orbit can be determined before the launch,
and the target may be either launched directly to that orbit or inserted into it after
launch, some time before the rendezvous manoeuvre.
The ascent trajectory of the interceptor must respect not only the requirements
imposed by the equations
GHAL GHA0 þ 2mp

tL ¼
xE

but also those imposed by lighting (see Sect. 4.3) and safety (see Sect. 4.4). In case
of impossibility of satisfying all of such requirements, the interceptor must have a
plane-change capability. By so doing, it is possible to rewrite the first of the two
equations shown above as follows
GHAL GHA0 þ 2m p
tL ¼ þ DTLW
xE
where DTLW indicates the width of the launch window, defined as shown below.
The plane-change capability is indicated by the angle Dη shown in the following
figure, due to the courtesy of NASA [45]. The change-of-plane manoeuvre is
assumed to be made at the rendezvous point defined by its argument of latitude, ut,
in the target orbit.
By the way, as has been shown in Sect. 1.9, the argument of latitude at epoch is
the angle, measured in the orbital plane, between the ascending node and the
position vector, r, of the satellite at epoch. Since the value of the angle Dη is small
in order to limit the propellant consumption, then the profile of the ascent trajectory
in case of a delayed launch is approximately the same as that in case of an on-time
launch. Therefore, the true anomaly, /s, of the interception point in the ascent
trajectory and the flight path angle, cs, have about the same values as those for a
nominal ascent trajectory.
The increment of velocity, Dv, required to change the orbital plane by Dη is then
h i12
Dv ¼ 2ðvs cos cs Þ2 ð1 cos D gÞ
as shown in the following figure.

The preceding equation, solved for Dη, yields

" 2 #
1 Dv
D g ¼ arccos 1
2 vs cos cs
According to Robertson [45], for a delayed launch, the inclination, is, of the
ascent trajectory plane with respect to the equator is
is ¼ arccosðcos Dg cos it sin Dg sin it cos ut Þ
where 0° < is < 180° for prograde orbits. The difference, DX, between the right
ascensions of the ascending nodes of the two ascent trajectories is
cos Dg cos it cos is

cos DX ¼
sin it sin is
sin Dg sin ut
sin DX ¼
sin is

sin DX
DX ¼ arctan
cos DX
The quantities Dkt and Dks are given by

tan uL
Dkt ¼ arcsin
tan it

tan uL
Dks ¼ arcsin
tan is
where
0
\Dkt;s \90
90
\Dkt;s \180
for southerly launches
and the notation Dkt,s indicates that these constraints hold for both Dkt and Dks.
Therefore, when the launch vehicle has a plane-change capability which
provides the necessary velocity increment Dv, then the delayed launch site may
rotate by an amount
DkLD ¼ DX þ Dks Dkt
past the on-time launch point.

Provided that the launch occurs within the launch window, it is also possible to
launch early instead of late, by an amount of time compatible with the plane-change
capability. In the case of an early launch, the inclination, is, of the ascent trajectory
plane with respect to the equator is
is ¼ arccosðcos Dg cos it þ sin Dg sin it cos ut Þ
where 0° < is < 180° for prograde orbits. Then, the remaining quantities (cos DX,
sin DX, DX, Dkt, and Dks) can be computed by using this value of is in the same
equations as those given above. Consequently, the early launch parameter, DkLE, is
DkLE ¼ DXE DkEs þ DkEt
The total launch window corresponding to the plane-change capability (Dv) of

the launch vehicle, that is, the time required for the Earth to rotate through the
angles DkLD and DkLE, is
DkL DkLD þ DkLE

¼
xE xE
where xE is the angular velocity at which the Earth rotates about its axis.
In case of a delayed launch, the width of the launch window is
DkLD
DTLW ¼
xE
In case of an early launch, the width of the launch window is
DkLE
DTLW ¼
xE
The two preceding expressions of DTLW (concerning the cases of, respectively,
late launch and early launch) yield the value of the term to be used on the right-hand
side of the equation given above
GHAL GHA0 þ 2mp

tL ¼ þ DTLW
xE
The equations shown above provide the basis for evaluating a possible direct
ascent trajectory along which an interceptor vehicle, launched from a site located on
the surface of the Earth, should move to approach at a very close distance a target
vehicle revolving around the Earth in an arbitrary orbit. Such equations are only
meant to define the fundamental requirements to be met by a direct ascent trajec-
tory. Further refinements are necessary for a precise rendezvous mission.
The second manner consists in allowing the interceptor to perform plane-change
manoeuvres. By so doing, the interceptor changes its orbital plane and reaches the
same orbital plane as the target. In practice, plane-change manoeuvres are expen-
sive, because they require a considerable amount of velocity increment (delta-v).
Therefore, the plane-change capability of the interceptor is usually limited to the
amount required to give some flexibility to the launch window during which a
direct ascent and rendezvous are possible.
This problem is greatly simplified by launching the target to a rendezvous-
compatible orbit, that is, to an orbit in which the target vehicle passes over the
launch site at least once per day. The perturbations to the central gravitational force
considered here are only those which are due to an oblate ellipsoid.
4.8 Rendezvous-Compatible Orbits
As has been anticipated in Sect. 4.7, the direct ascent of a space vehicle and its
rendezvous with another space vehicle become simpler when the target vehicle is
placed in a rendezvous-compatible orbit, that is, in an orbit in which the target
vehicle passes over the launch site at least once per day.
In other words, according to Perry [54], the target vehicle must be at the latitude
of the initial injection point over the launch site just as this point passes through the
orbital plane. Since a revolution of the target vehicle with respect to a latitude circle
requires a nodal period T♌ (see below) of the same vehicle to be accomplished, then
there must be an integral number of nodal periods between consecutive passages of
the sub-injection point through the same portion of the orbital plane, as shown by
the following figure, due to the courtesy of NASA [54].
According to Robertson [45], in order for the target vehicle to pass over the
launch site every nth revolution after launch, the following relation must hold:
where n is an integral number of revolutions, T♌ is the nodal period of the target

vehicle over an oblate Earth (that is, the time which elapses between two successive
south–north passages of the target vehicle through the equatorial plane of the
Earth), m is an integral number of days between passages, X′t is the nodal
regression per revolution, tascent is the time which elapses from the moment in
which the target vehicle is launched to the moment in which it is inserted into orbit,
and xE is the angular velocity of rotation of the Earth.
As has been shown in Sect. 4.7, the regression (for a prograde orbit) rate of the
right ascension of the ascending node is
2
3 rE cos it
X0t ¼ nt C20 2
2 at 1 e2 t
where rE is the equatorial radius of the Earth, at, et, and nt are, respectively, the
major semi-axis, the eccentricity, and the mean motion of the orbit of the target
vehicle, and finally C20 = −0.001082636 is the principal zonal coefficient of the
gravitational field of the Earth. Therefore, the nodal regression in one revolution
(2p rad) results from the preceding expression multiplied by 2p, which yields
2
rE cos it
X0t ¼ 3p nt C 20 2
at 1 e2 t
4.8 Rendezvous-Compatible Orbits 649
The nodal period of the target vehicle is computed by Robertson [45] by means
of the following expression due to Kalil and Martikan [55]:
The expression written above holds when et is small, that is, when et is of the
order of magnitude of C20. The computation of T♌ for the general case has been
considered by several authors, as will be shown at the end of this paragraph.
The previous expression
solved for T♌ yields the required nodal period for a rendezvous-compatible orbit
such that the vehicle to be intercepted passes above the launch site at least once per
day (m = 1) every n revolutions
The terms relating to inclination, it, are included into the expressions of X′t and
T♌; however, as Robertson [45] points out, the equation written above does not
depend on the inclination of the rendezvous-compatible orbit to the first order.
When it is desired to take advantage of the two launch opportunities per day, then it
is necessary to specify an inclination. In other words, when it is desired to make a
southerly pass over the launch site q revolutions after the northerly pass in a period
of 24 h, then the inclination of the rendezvous-compatible orbit results from the
following equations, in which a circular rendezvous-compatible orbit is assumed.
These equations, indicated by Robertson [45], are given below.
where
Du ¼ 2 arctan ðcot u L cos AÞ

DkE ¼ 2 arccot ðtan A sin u L Þ

cos it
A ¼ arcsin
cos u L
In addition, u*L and A are, respectively, the latitude of the launch site and the
launch azimuth, which falls between 0° and 90° when the launch site is in the
northern hemisphere of the Earth, as shown in the following figure, due to the
The equations given above, that is,
must be solved numerically, because X′t and T♌ depend on the inclination, it, of the
orbit of the target vehicle with respect to the equatorial plane of the Earth.
According to Robertson [45], a first-trial value of T♌, to be introduced into the
equation T♌ = (2p + X′t n − tascent xE)/(nxE), results from T♌ 2p(a3t /lE)½,
which means that, in the first trial of the trial-and-error procedure, the oblate shape
of the Earth is neglected, and therefore the nodal period of the target vehicle is
approximated by its Keplerian period T = 2p(a3t /lE)½.
The following section of this paragraph concerns the question of the nodal
period of an artificial satellite. It is not our intention to add anything to what has
been written by several authors on this subject. We limit ourselves to give a brief
account of the matter to the reader.
King-Hele [56, page 63, Eq. 73] defines the nodal period of a satellite as the time
which elapses between two successive northward crossings of the equator by the
4.8 Rendezvous-Compatible Orbits 651
satellite. He computes the nodal period, up to and including the second-order

gravitational perturbations, as follows
which differs slightly from the expression given by Kalil and Martikan [55].
El’yasberg [57, page 217, Eq. 13.41] defines the nodal period of a satellite as
follows
where t is the flight time and u is the angular distance from the current point to the
node of the corresponding osculating orbit.
Merson defines the nodal period, T♌, as the change in time, measured from the
ascending node, as the argument of latitude (u = x + /) goes from 0 to 2p.
Merson [58, page 35, Eq. 75] computes T♌ as follows
where p = a(1 − e2) and n0 = (lE/a3)½ are, respectively, the semi-latus rectum and
the mean motion at the node of the osculating orbit.
Blitzer [59, page 1459, Eq. 1], or Ref. [60, page 42, Eq. 5.25], defines T♌ as the
time interval from one to the succeeding ascending node, and computes T♌ as
follows
Thomas [61, page 632, Eq. 99] computes T♌ as follows

Rees [62, page 330, Eq. 10.13] defines T♌ as the time between successive
ascending or descending nodes, and computes T♌ as follows
Other authors than those named above, for example, Vallado [63], Tapley et al.
[64], Mortari et al. [65], and Noreen et al. [66], define, all of them, the nodal period
as follows
where n is the unperturbed mean motion, and x′ and M′0 are the rates of change in
time of, respectively, the argument of perigee and the mean anomaly due to the
perturbing forces. With this definition, as is the case with the former one, the nodal
period T♌ coincides with the Keplerian period T = 2p/n in the absence of pertur-
bations (x′ = M′0 = 0).
By introducing the secular rates of change in time of x′ and M′ due to the oblate
shape of the Earth (see Sect. 3.2), that is,
3 r 2 1 5 cos2 i
x0 ¼ nC 20
E
4 a ð 1 e2 Þ 2
3 r 2 1 3 cos2 i
M 0 ¼ n þ nC20
E
3
4 a ð 1 e2 Þ 2
into the latter definition of T♌, we obtain
where, again, T = 2p/n is the Keplerian period of revolution.

4.9 Intermediate Orbits for Rendezvous 653
4.9 Intermediate Orbits for Rendezvous
The present paragraph deals with the techniques used for inserting the intercepting
vehicle into an intermediate parking orbit before its final transfer and rendezvous
with the target vehicle. By so doing, it is possible to launch the intercepting vehicle
at nearly any time in which the launch site passes through the plane containing the
orbit of the target vehicle. The proper phasing conditions, which are necessary for a
rendezvous between the two vehicles, are achieved subsequently, when the inter-
cepting vehicle has already been inserted into its intermediate orbit. By phasing
conditions, we mean the proper position of the intercepting vehicle as a function of
time which is necessary for a successful rendezvous manoeuvre. In other words,
this manoeuvre must occur when the intercepting vehicle has reached an appro-
priate position along its orbit, because the interceptor and the target must arrive at
the rendezvous point at the same time.
In most cases, the intermediate orbit in which the interceptor is to be placed and
the orbit of the target are both of them circular, and the rendezvous manoeuvre is a
motor firing which injects the interceptor from an appropriate position along its
initial circular orbit into a transfer ellipse which intersects the target orbit at the
rendezvous point. In line of principle, the same holds in case of a transfer from one
to another elliptical orbit, but in this case the problem must be solved numerically.
The first technique discussed here considers the case in which both the inter-
ceptor (also called the shuttle) and the target move along circular orbits about the
Earth, as shown in the following figure, which is due to the courtesy of NASA [45].
In the case illustrated below, the injection of the shuttle into the transfer ellipse
occurs at the perigee of this ellipse. In addition, the flight path angle, cR, of the
transfer ellipse at the point of intersection with the target orbit has a certain value
which is not necessarily equal to zero. In the particular case of cR = 0, the transfer
ellipse is tangent to the circular orbit of the target at the rendezvous point, and the
transfer ellipse is a Hohmann-type trajectory, which has been shown in Sect. 4.6.
Robertson [45] supposes the value of the angle cR to be chosen by the designer
and computes the elements of the transfer ellipse as a function of cR.
We prefer to do the same thing by using another method, suggested by Braeunig
[67], which is shown below. Let aTR be the major semi-axis of the transfer ellipse.
With reference to the preceding figure, we choose the value of aTR such that
1
aTR [ ðRI þ RT Þ
2
where RI and RT are the radii of the two circular orbits. As a result of this choice, the
transfer ellipse intersects the circular orbit of the target vehicle in the rendezvous
point at some flight path angle cR 6¼ 0. Since the interceptor is injected into the
transfer ellipse at the perigee (of radius rP = RI) of this ellipse, then the eccentricity,
eTR, of the transfer ellipse results from
RI
eTR ¼ 1
aTR
This is because the equation a = ea + rP holds with any elliptic orbit.

The semi-latus rectum, pTR, of the transfer ellipse results from
pTR ¼ aTR ð1 e2TR Þ
With reference to the preceding figure, since the equation r = p/(1 + e cos /)
holds with any elliptic orbit, then we substitute r = RT, p = pTR, e = eTR, and
/ = /R, and solve the preceding equation for cos /R. This yields

pTR
1
r2
cos /R ¼
eTR
Hence, we compute /R by means of the inverse cosine function. The value

chosen for aTR is such as to assure us that /R is a few degrees less than 180°.
Since the flight path angle, c, at any point of an elliptic orbit is by definition
vr r0 e sin /
tan c ¼ ¼ 0¼
v/ r/ 1 þ e cos /
where vr and v/ are, respectively, the radial component and the transverse com-
ponent of the velocity vector, v, of the interceptor in any point of the transfer
ellipse, then we substitute c = cR, e = eTR, and / = /R, and compute cR.
By using the following expressions of Sect. 1.3:
where e = eTR and / = /R, the eccentric anomaly, ÆR, at the rendezvous point
results from
Likewise, by using the following expression of Sect. 1.3:
where e = eTR and Æ = ÆR, it is possible to compute the mean anomaly, MR, at the
rendezvous point.
The transfer time, DtR, of the interceptor along the transfer ellipse, from the point
of perigee to the point of rendezvous, results from
12
a3TR
DtR ¼ MR
lE
With reference to the preceding figure, the lead angle, w, measured from the
radius vector of the interceptor to the radius vector of the target at the moment in
which the interceptor is injected into the transfer ellipse, results from
w ¼ /R nT DtR
where nT = (lE/R3T)½ is the mean motion of the circular orbit of the target vehicle.
The coasting time, tcoasting, spent by the interceptor in the intermediate orbit
results from
w w0
tcoasting ¼
nT nI
where, with reference to the preceding figure, w0 is the lead angle measured from
the radius vector of the interceptor to the radius vector of the target at the moment in
which the interceptor is injected into the intermediate circular orbit, and the mean
motion nI = (lE/R3I )½ is the mean motion of the intermediate circular orbit.
The difference of velocity (DvTR), which is needed at the moment in which the
interceptor is inserted from its circular intermediate orbit to its elliptic transfer orbit,
results from the vis-viva integral v2/l = 2/r − 1/a, and is therefore
12 12
2 1 l
DvTR ¼ lE E
RI aTR RI
The previous equation holds because the two velocity vectors are aligned.
At the point of rendezvous, the magnitude, vR, of the velocity vector vR tangent
to the transfer ellipse is
12
2 1
vR ¼ lE
RT aTR
and the magnitude, vT, of the velocity vector vT tangent to the circular orbit of the
target vehicle is
12
lE
vT ¼
RT
Since cR is the angle between the two velocity vectors named above, then the
difference of velocity (DvR), which is needed at the moment in which the interceptor
is inserted from its elliptic transfer orbit to the circular orbit of the target vehicle,
results from the cosine rule
1
DvR ¼ v2R þ v2T 2vR vT cos cR 2
The total difference of velocity needed for the rendezvous manoeuvre is

therefore
Dvtotal ¼ DvTR þ DvR
If the value of the major semi-axis, aTR, of the transfer ellipse were chosen equal
to ½(RI + RT), then the transfer ellipse would be a Hohmann-type ellipse, and the
value of the flight path angle at the rendezvous point would be equal to zero, as has
been shown in Sect. 4.6.
The second technique discussed here is called the co-orbital technique, and uses
an intermediate, or phasing, elliptic orbit, which is shown in the following figure,
due to the courtesy of NASA [45]. The phasing orbit intersects the circular orbit of
the target vehicle or is tangent to this orbit, such that the interceptor be at a given
time either slightly ahead or behind the target, but that, after a number m of rev-
olutions, this difference of position between the interceptor and the target should
reduce to zero.
The co-orbital technique has the advantage of simplifying the phasing

manoeuvre, because the interceptor must only be inserted into an elliptic orbit
whose period be such as to bring the two vehicles to the same position after
m revolutions, but also the disadvantage of requiring a small difference of position
between the two vehicles, in order to keep the number m of revolutions small. As
shown in the preceding figure, the interceptor is inserted into a phasing elliptic orbit
whose period makes it possible to the target to reach the perigee of the ellipse after a
small number of revolutions of the interceptor in the same orbit. In the figure, the
target is supposed to stay behind the interceptor by an angle Dw in the moment in
which the interceptor is inserted into the elliptic phasing orbit.
The difference of time existing between the two vehicles in this moment is
Dw
Dt ¼
nT
where nT = (lE/a3T)½ is the mean motion of the target vehicle.

In order for a rendezvous between the two vehicles to occur after m revolutions
of the interceptor in its elliptic phasing orbit, the period TPH of this orbit must be
chosen so that
Dt
TPH ¼ TT þ
m
where TT = 2p(R3I /lE)½ is the period of the target in its circular orbit of radius RI.
After determining the period, TPH, of the interceptor in its phasing orbit as
indicated above, the major semi-axis, aPH, of this elliptic orbit results from solving
the equation TPH = 2p(a3PH/lE)½ for aPH. This yields
" 2 #13
TPH
aPH ¼ lE
2p
The difference of velocity (Dv), needed at the moment in which the interceptor
reaches the target at the point of rendezvous, results from the vis-viva integral
v2/l = 2/r − 1/a, and is therefore
12 12
2 1 l
Dv ¼ lE E
RI aPH RI
The time required for rendezvous is
trendezvous ¼ mTPH
The third technique discussed here, called the line-of-sight delta-v, is a particular
case of the first technique discussed previously.
The third technique differs from the first, because in this case the rocket motor is
fired along the line of sight between the interceptor and the target. With reference to
the following figure, due to the courtesy of NASA [45], let RI be the radius of the
circular intermediate orbit into which the interceptor is inserted.
According to this technique, the axis of the rocket motor is aligned to form a
predetermined angle, d, with respect to the local horizontal (that is, with respect to
the line perpendicular to RI). When the line of sight between the interceptor and the
target forms the same angle, d, with respect to the local horizontal, then the rocket
motor of the interceptor is fired. As a result of the thrust applied, the interceptor is
inserted into an elliptic transfer trajectory which keeps the two vehicles along the
line of sight from the moment of motor firing to the moment of rendezvous.
The value of the angle d is chosen by the designer. The corresponding value of the
increment of velocity DvTR, which is to be applied to the interceptor at the moment
of its insertion into the transfer ellipse, is to be found by iteration on the following
equations.
Let Dvest be an estimated value of the unknown impulse DvTR. With reference to
the preceding figure, before the application of the velocity impulse Dvest, the
velocity vector of the intercepting vehicle has only the transverse (or circumfer-
ential or horizontal) component vI. When the velocity impulse Dvest is applied, the
velocity vector of the intercepting vehicle has two components:
• the radial component Dvest sin d; and
• the transverse component vI + Dvest cos d.
The magnitude of the vector vTR, resulting from the sum of these two compo-
nents, is
h i12 1
vTR ¼ Dv2est sin2 d þ ðvI þ Dvest cos dÞ2 ¼ Dv2est þ v2I þ 2D vest vI cos d 2
The vector vTR is the velocity of the interceptor at the time of its insertion into
the transfer ellipse. The position vector of the interceptor, RI, at the same time has
magnitude RI and is directed radially.
The components of these two vectors along the directions, respectively, radial,
transverse, and bi-normal are the following
2 3 2 3
RI Dvest sin d
RI 4 0 5 vTR 4 vI þ Dvest cos d 5
0 0
Consequently, the components of the moment of momentum per unit mass
hI ¼ RI vTR
along the same directions, at the time of the interceptor insertion into the transfer
ellipse, are the following
2 3
0
hI 4 0 5
RI ðvI þ Dvest cos dÞ
By definition, the flight path angle c of a spacecraft at any time results from
h ¼ rv cos c
Therefore, there results at any time

1
12
h2 2
sin c ¼ 1 cos c ¼ 1 2 2
2
r v
At the time of the interceptor insertion into the transfer ellipse, the position, the
velocity, and the moment of momentum per unit mass of the interceptor have the
components given above. By substituting the values of these components in the
preceding equation, there results, after simplification, the value of flight path angle
of the interceptor at the time of its insertion into the transfer ellipse

Dvest sin d
cTR ¼ arcsin
vTR
The magnitude of the moment of momentum per unit mass of the interceptor
along its transfer ellipse results from
hTR ¼ RI vTR cos cTR
Remembering (see Sect. 1.1) the equality p = h2/lE, the semi-latus rectum of the
transfer ellipse results from
h2TR
pTR ¼
lE
By using the vis-viva integral (v2/lE = 2/r − 1/a) and solving for a, the major
semi-axis of the transfer ellipse results from
lE RI
aTR ¼
2lE v2TR RI
Remembering that p = a(1 − e2) and solving for e, the eccentricity of the
transfer ellipse results from
12
pTR
eTR ¼ 1
aTR
Now we use the same equations as those shown above for the first technique to
compute the true anomaly, /TR, of the interceptor at the time of its insertion into the
transfer ellipse. These equations are the following
pTR RI
cos /TR ¼
RI eTR
pTR tan cTR
sin /TR ¼
RI eTR

sin /TR
/TR ¼ arctan
cos /TR
By operating likewise, the eccentric anomaly and the mean anomaly of the
interceptor at the time of its insertion into the transfer ellipse result from
At the point of rendezvous, if the quantities computed above had the correct
values, then the velocity, vR, of the interceptor in its transfer orbit would be given
by the vis-viva integral (v2/lE = 2/r − 1/a), as follows
12
2 1
vR ¼ lE
RT aTR
where RT is the radius of the circular orbit of the target vehicle.

Remembering the equation h = rv cosc, this value of vR makes it possible to
compute the flight path angle, cR, of the interceptor at the rendezvous point as
follows

hTR
cR ¼ arccos
RT vR
By using this value of cR, it is possible to compute the true anomaly of the
interceptor at the point of rendezvous as follows
pTR RT
cos /R ¼
RT eTR
pTR tan cR
sin /R ¼
RT eTR

sin /R
/R ¼ arctan
cos /R
Likewise, the eccentric anomaly, ÆR, and the mean anomaly, MR, of the
interceptor at the point of rendezvous result from
The time, which elapses from the time of insertion of the interceptor into the
transfer ellipse to the time of rendezvous, results from
3 12
aTR
DtR ¼ ðMR MTR Þ
lE
During this time, the angle between the radius vector relating to the insertion of
the interceptor into the transfer ellipse and the radius vector relating to the ren-
dezvous point of the transfer ellipse is D/R = /R − /TR.
However, the preceding figure, due to the courtesy of NASA [45], shows that the
angle through which the target vehicle must travel along its circular orbit during the
time DtR in order for a rendezvous to occur is
D/T ¼ D/R b
where the angle b is to be computed as will be shown below.

Let us consider again the preceding figure. By applying the sine rule (the sides of
any triangle are proportional to the sines of the opposite angles), there results
sinðp=2 þ dÞ sin a
¼
RT RI
Since sin(p/2 + d) = cosd, then the preceding equation, solved for a, yields

RI cos d
a ¼ arcsin
RT
Hence
p
b¼ da
2
So far, the motion of the target vehicle has not been considered. According to the
hypotheses made at the beginning of this paragraph, the two vehicles revolve about
the Earth in circular orbits. Therefore, the arc of circumference travelled by the
target vehicle during the time interval, DtR, computed above is
DtR nT
With reference to the preceding figure, the phase angle error, at the time in which
the intercepting vehicle reaches the orbit of the target vehicle, is
D/phase ¼ D/R DtR nT
Now, it is necessary to perform a further iteration, by using a new value of Dvest,

for the purpose to reducing the phase angle error to an acceptable value.
At the end of this iterative process, when the correct value of Dvest has been
determined, we set
DvTR ¼ Dvest
The equations shown above (with DvTR instead of Dvest) can be used to compute
the elements of the transfer orbit. These equations are rewritten below for
convenience of the reader. The velocity of the interceptor at the time of its insertion
into the transfer ellipse is
1
vTR ¼ Dv2TR þ v2I þ 2D vTR vI cos d 2
where, again, the value of the angle d is chosen by the designer.

The flight path angle of the interceptor at the time of its insertion into the transfer
ellipse is

D vTR sin d
cTR ¼ arcsin
vTR
The moment of momentum per unit mass of the interceptor along its transfer
ellipse is
hTR ¼ RI vTR cos cTR
The semi-latus rectum of the transfer ellipse is
h2TR
pTR ¼
lE
The major semi-axis of the transfer ellipse is
lE RI
aTR ¼
2lE v2TR RI
The eccentricity of the transfer ellipse is

12
pTR
eTR ¼ 1
aTR
The true anomaly of the interceptor at the time of its insertion into the transfer
ellipse is
pTR RI
cos /TR ¼
RI eTR
pTR tan cTR
sin /TR ¼
RI eTR

sin /TR
/TR ¼ arctan
cos /TR
The eccentric anomaly and the mean anomaly of the interceptor at the time of its
insertion into the transfer ellipse are
The velocity of the interceptor in its transfer orbit at the point of rendezvous is
12
2 1
vR ¼ lE
RT aTR
The flight path angle of the interceptor at the point of rendezvous is

hTR
cR ¼ arccos
RT vR
The true anomaly of the interceptor at the point of rendezvous is
pTR RT
cos /R ¼
RT eTR
pTR tan cR
sin /R ¼
RT eTR

sin /R
/R ¼ arctan
cos /R
The eccentric anomaly and the mean anomaly of the interceptor at the point of
rendezvous are
The time of flight of the interceptor along the transfer ellipse is

3 12
aTR
DtR ¼ ðMR MTR Þ
lE
The lead angle, w, measured from the position of the interceptor to the position of
the target at the moment in which the interceptor is injected into the transfer ellipse, is
w ¼ /R nT DtR
The coasting time tcoasting, spent by the interceptor in the intermediate orbit, is
w w0
tcoasting ¼
nT nI
where w0 is the lead angle measured from the position of the interceptor to the
position of the target at the moment in which the interceptor is injected into the
intermediate circular orbit, and nI = (lE/R3I )½ is the mean motion of the interme-
diate circular orbit.
By operating as has been shown above for the first of these three techniques, at
the point of rendezvous, the magnitude, vR, of the velocity vector vR tangent to the
transfer ellipse is
12
2 1
vR ¼ lE
RT aTR
and the magnitude, vT, of the velocity vector vT tangent to the circular orbit of the
target vehicle is
12
lE
vT ¼
RT
Since cR is the angle between the two vectors named above, then the difference of
velocity (DvR), needed at the moment in which the interceptor is inserted from its elliptic
transfer orbit into the circular orbit of the target vehicle, results from the cosine rule
1
DvR ¼ v2R þ v2T 2vR vT cos cR 2
The total difference of velocity required for the rendezvous manoeuvre is
Dvtotal ¼ DvTR þ DvR

4.10 The Hill–Clohessy–Wiltshire Equations
The equations shown below, so called after Hill [68] and Clohessy and Wiltshire
[69], describe the motion of one of two vehicles involved in a rendezvous
manoeuvre with respect to a system of reference, Oxyz, whose origin, O, is in the
centre of mass of the other vehicle. With reference to the following figure, let us
consider the system of reference Oxyz, whose origin, O, is in the centre of mass of
the target vehicle, which in turn is assumed to revolve about the centre of mass, G,
of the Earth in a circular orbit of radius R.
The axes of this system of reference are such that the x-axis (whose unit vector is
ux) is the radial direction, positive for increasing values of the radius vector R, the y-
axis (whose unit vector is uy) is the transverse direction, positive for increasing
values of the true anomaly /, and the z-axis (whose unit vector is uz) is perpen-
dicular to the plane of motion, such that uz = ux uy. Of course, the system of
reference Oxyz is not inertial, and rotates with respect to a hypothetic inertial system
of reference having its origin in the centre of mass of the Earth with an angular
velocity vector x = xuz, whose magnitude depends on the radial distance, R,
between the target vehicle and the centre of mass of the Earth, that is, x = (lE/R3)½,
where lE is the gravitational parameter of the Earth.
Let us consider the other vehicle, which moves in the vicinity of the origin O.
Let R and r be the geocentric position vectors of, respectively, the target and the
interceptor. Let q be the position vector of the interceptor with respect to the system
of reference Oxyz, such that the magnitude q of the vector q is assumed to be much
smaller than R. This assumption is reasonable, since the distance between the target
and the interceptor is much smaller than the radius of the circular orbit of the target.
We want to introduce the vector equation
r ¼ Rþq
4.10 The Hil–Clohessy–Wiltshire Equations 669
into the equation of motion of the target, that is,

lE
r00 þ r¼a
r3
Then we want to take advantage of the property q R to expand r−3 in a

Maclaurin series truncated after the first order. This makes it possible to approxi-
mate the preceding differential equation, which is not linear, to a linear differential
equation.
As is well known, the preceding differential equation holds when the position
vector, r, and the perturbing acceleration vector, a, of an object are related to an
inertial system of reference. By contrast, when the motion of that object is related to
a non-inertial reference system, as is the case with the Oxyz system, then the
acceleration vector, a, must include not only the perturbing accelerations but also
the fictitious accelerations.
By writing R + q in place of r into the differential equation written above, there
results
lE
R00 þ q00 ¼ ð R þ qÞ þ a
jR þ qj3
The denominator of the fraction on the right-hand side of the equation written
above is a scalar quantity. The reciprocal of this quantity may be written as follows
jR þ qj3 ¼ ½ðR þ qÞ ðR þ qÞ2 ¼ ½ðR RÞ þ 2ðR qÞ þ ðq qÞ2

3 3
3
R q q q 2
¼ R3 1 þ 2 2 þ 2
R R
By setting for convenience 2(R q)/R2 + (q q)/R2 x and remembering the

following expansion of (1 + x)p in a Maclaurin series, there results
pðp 1Þ 2 pðp 1Þðp 2Þ 3

ð1 þ xÞp ¼ 1 þ px þ x þ x þ
2! 3!
By setting p = −3/2 and expanding up to the first power of x, there results
3
ð 1 þ xÞ 2 ¼ 1 x þ O x2
3
2
and therefore
3
R q q q 2 3 Rq qq
jR þ qj3 ¼ R3 1 þ 2 2 þ 2 R3 1 2 2 þ 2
R R 2 R R

Rq
R3 1 3 2
R
The fictitious accelerations, aF, which are due to the non-uniform motion of the
non-inertial reference system Oxyz with respect to an inertial reference system, are
aF ¼ 2x q0 x ðx qÞ x0 q
On the right-hand side of the preceding equation, the first term is the Coriolis
acceleration, the second is the centrifugal acceleration, and the third is the Euler
acceleration. In the present case, the Euler acceleration is zero, because the target
vehicle has been assumed to revolve about the Earth in a circular orbit at a uniform
angular velocity x. Let aP be the resultant of all the perturbing accelerations. When
the equation of motion r″ + (lE/r3)r = a is written with respect to the non-inertial
reference system Oxyz, the acceleration, a, to be taken into account on the
right-hand side of this equation is
a ¼ aF þ a P
where aF has the expression written above. Therefore, the equation of motion of an
object with respect to the non-inertial reference system Oxyz is

lE 00 0 lE 3
3 R þ q þ 2x q þ x ðx qÞ ¼ 3 R þ q 2 ðR qÞR þ aP
R R R
where the first term on the left-hand side of the preceding equation is the inertial
acceleration vector R″. The preceding equation becomes after simplification

lE 3
q00 þ 2x q0 þ x ðx qÞ ¼ q ð R qÞR þ aP
R3 R2
The vectors appearing in this equation have the following components along the
axes x, y, and z:
R ¼ Rux
q ¼ xux þ yuy þ zuz
x ¼ xuz
aP ¼ aPx ux þ aPy uy þ aPz uz
Therefore, the preceding equation of motion, projected onto the axes x, y, and
z defined above, yields the following three scalar equations:
x00 2x y0 3x2 x ¼ aPx

y00 þ 2x x0 ¼ aPy
z00 þ x2 z ¼ aPz
which are the Hill–Clohessy–Wiltshire equations. Some authors, Shepperd [70] for
one, make a different choice for the axes x, y, and z.
The coefficients on the left-hand side are constant. The two equations along
x and y are coupled, whereas the equation along z is not coupled with the first two.
The conditions at t = 0 are x = x0, y = y0, z = z0, x′ = x′0, y′ = y′0, and z′ = z′0. In
other words, the three components of the position vector q of the interceptor and the
three components of the velocity vector q′ of the interceptor must be given at t = 0.
These equations can be solved analytically, as will be shown below.
Let us consider the corresponding homogeneous differential equations, that is,
x00 2x y0 3x2 x ¼ 0
y00 þ 2x x0 ¼ 0
z00 þ x2 z ¼ 0
relating to the case of absence of perturbations. The second of these equations may
0
ðy0 þ 2x xÞ ¼ 0
that is,
y0 þ 2x x ¼ constant
The constant on the right-hand side is determined by imposing the following

initial condition
y0 þ 2x x ¼ y00 þ 2x x0
which may also be written as follows
y0 ¼ y00 þ 2x ðx0 xÞ
By substituting this value of y′ into the first (x″ − 2xy′ − 3x2x = 0) of the three
homogeneous differential equations, there results

x00 2x y00 þ 2x ðx0 xÞ 3x2 x ¼ 0
which yields after simplification
x00 þ x2 x ¼ 2x y00 þ 4x2 x0
where the quantity on the right-hand side of the preceding equation is constant with
time. By setting for convenience C = 2xy′0 + 4x2x0, there results
x00 þ x2 x ¼ C ¼ constant
As is well known, the general solution of the preceding non-homogeneous

differential equation results from summing the general solution of the corre-
sponding homogeneous differential equation (x″ + x2x = 0) with a particular
solution, x*, of the non-homogeneous differential equation (x″ + x2x = C).
Let us find first a particular solution of x″ + x2x = C. Since C is constant, there
results evidently
C 2x y00 þ 4 x2 x0
x ¼ ¼
x2 x2
The general solution of x″ + x2x = 0 is evidently x = A sin(xt) + B cos(xt),

where A and B are constants to be determined; therefore, the general solution of the
non-homogeneous equation x″ + x2x = 2xy′0 + 4x2x0 is
2x y00 þ 4x2 x0
x ¼ A sinðx tÞ þ B cosðx tÞ þ
x2
The values of the constants A and B result from imposing the initial conditions,
as will be shown below.
First condition:
x ð 0Þ ¼ x 0
Hence
2x y00 þ 4x2 x0
Bþ ¼ x0
x2
which, solved for B, yields
y00
B ¼ 3x0 2
x
Second condition:
x0 ð0Þ ¼ x00
Hence Ax = x′0, which, solved for A, yields
x00
A¼
x
These values of A and B, put into x = A sin(xt) + B cos(xt) + (2xy′0 + 4x2x0)/x2,

yield

x00 y0 y0
x¼ sinðx tÞ 3x0 þ 2 0 cosð x tÞ þ 2 0 þ 4x0
x x x

1 2
x ¼ ½4 3 cosðx tÞx0 þ sinðx tÞ x00 þ ½1 cosðx tÞy00
x x
The preceding equation, differentiated with respect to time, yields
x0 ¼ ½3 x sinðx tÞx0 þ ½cosðx tÞx00 þ ½2 sinðx tÞy00
The equation x = [4 – 3 cos(xt)]x0 + (x′0/x)sin(xt) + (2y′0/x)[1 − cos(xt)],

substituted into y′ = y′0 + 2x(x0 − x), yields

0 x00 y00
y ¼ y00 þ 2 x x0 ½4 3 cosðx tÞx0 sinðx tÞ 2 ½1 cosðx tÞ
x x
which in turn, after simplification, yields
y0 ¼ 6x½cosðx tÞ 1x0 2x00 sinðx tÞ þ ½4 cosðx tÞ 3y00
The preceding expression, integrated with respect to time, yields

1 2x0 4
y ¼ 6x sinðx tÞ t x0 þ 0 cosðx tÞ þ sinðx tÞ 3t y00 þ C
x x x
The preceding expression, evaluated at t = 0, yields
2x00
y0 ¼ þC
x
which in turn, solved for C, yields
2x00
C ¼ y0
x
This value of C, substituted into

1 2x00 4
y ¼ 6x sinðx tÞ t x0 þ cosðx tÞ þ sinðx tÞ 3t y00 þ C
x x x
yields

2x00 4
y ¼ 6½sinðxtÞ xtx0 þ ½cosðxtÞ 1 þ sinðxtÞ 3t y00 þ y0
x x
Let us consider now the third differential equation z″ + x2z = 0. As is well

known, the general solution of this equation is
z ¼ A sinðxtÞ þ B cosðxtÞ
The values of the constants A and B result from imposing the initial conditions,
as will be shown below. The preceding expression, differentiated with respect to
time, yields
z0 ¼ A x cosðxtÞ B x sinðxtÞ
Imposing the initial conditions in the two preceding expressions yields
z0 ¼ B
z00 ¼ Ax
Hence, the general solution of the differential equation z″ + x2z = 0 is

1
z ¼ ½cosðx tÞz0 þ sinðx tÞ z00
x
z0 ¼ ½x sinðx tÞz0 þ ½cosðx tÞz00

Summarising, the expressions derived above

1 2
x ¼ ½4 3 cosðx tÞx0 þ sinðx tÞ x00 þ ½1 cosðx tÞy00
x x

2 4
y ¼ 6½sinðx tÞ x tx0 þ y0 þ ½cosðx tÞ 1x00 þ sinðx tÞ 3t y00
x x

1
z ¼ ½cosðx tÞz0 þ sinðx tÞ z00
x
x0 ¼ ½3x sin ðx tÞ x0 þ ½cosðx tÞx00 þ ½2 sin ðx tÞy00
y0 ¼ 6x½cosðx tÞ 1x0 þ ½2 sinðx tÞx00 þ ½4 cosðx tÞ 3y00
z0 ¼ ½x sinðx tÞz0 þ ½cosðx tÞz00
are the analytical solution of the Hill–Clohessy–Wiltshire equations. By using the

matrix notation, this solution may be written as follows
qðtÞ ¼ Urr ðtÞq0 þ Urv ðtÞq00

q0 ðtÞ ¼ Uvr ðtÞq0 þ Uvv ðtÞq00
where q(t) and q′(t) are, respectively, the position vector and the velocity vector of
an object with respect to the non-inertial reference system Oxyz, and the four 3 3
matrices Urr(t), Urv(t), Uvr(t) and Uvv(t) have the following elements
2 3
4 3 cosðx tÞ 0 0
6 7
Urr ðtÞ 4 6½sinðx tÞ x t 1 0 5
0 0 cosðx tÞ
2 3
ð1=xÞ sinðx tÞ ð2=xÞ½1 cosðx tÞ 0
6 7
Urv ðtÞ 4 ð2=xÞ½cosðx tÞ 1 ð4=xÞ sinðx tÞ 3t 0 5
0 0 ð1=xÞ sinðx tÞ
2 3
3 x sinðx tÞ 0 0
6 7
Uvr ðtÞ 4 6 x½cosðx tÞ 1 0 0 5
0 0 x sinðx tÞ
2 3
cosðx tÞ 2 sinðx tÞ 0
6 7
Uvv ðtÞ 4 2 sinðx tÞ 4 cosðx tÞ 3 0 5
0 0 cosðx tÞ
Let s (a function of time) and s0 (a constant) be two 6 1 column vectors whose

components are
2 3 2 3
x x0
6y7 6 y0 7
6 7 6 7
6z7 6 z0 7
s6 7
6 x0 7 s0 6
6 x0 7
7
6 07 6 00 7
4y 5 4y 5
0
z0 z00
The analytical solution of the Hill–Clohessy–Wiltshire equations may be briefly

written as follows
s ¼ UðtÞs0
where the 6 6 matrix U(t), which is called the state transition matrix for the
Hill–Clohessy–Wiltshire equations from epoch 0 to epoch t, is given below.

4 3 cos xt 0 0 ð1=xÞ sin xt ð2=xÞð1 cos xtÞ 0

6ðsin xt xtÞ ð2=xÞðcos xt 1Þ ð4=xÞ sin xt 3t
1 0 0

0 0 cos xt 0 0 ð1=xÞ sin xt

3x sin xt
0 0 cos xt 2 sin xt 0

6xðcos xt 1Þ 0 0 2 sin xt 4 cos xt 3 0

0 0 x sin xt 0 0 cos xt
As is easy to verify, the four matrices Urr(t), Urv(t), Uvr(t), and Uvv(t) are
partitions of the state transition matrix U(t).
4.11 The Hill–Clohessy–Wiltshire Equations Applied

to Rendezvous Manoeuvres
The case considered below is a rendezvous manoeuvre, which comprises two

impulsive firings of the rocket motor of an intercepting vehicle directed towards a
target vehicle. The first of these firings is performed at the initial time t = 0, when
the interceptor is injected from its parking orbit around the Earth into a transfer
ellipse which intersects the circular orbit of the target vehicle at the point of ren-
dezvous. The second firing is performed at the final time t = tf, when the interceptor
leaves the transfer ellipse at the point of rendezvous in order to acquire not only the
position but also the velocity of the target vehicle.
Let 0− and 0+ be the times, respectively, immediately before and immediately
after the time t = 0 at which the first firing of the rocket motor is performed.
Likewise, let t−f and t+f be the times, respectively, immediately before and imme-
diately after the time tf at which the second firing of the rocket motor is performed.
4.11 The Hill-Clohessy-Wiltshire Equations Applied to Rendezvous … 677
We want to determine the components x′+0 , y′+0 , and z′+0 of the velocity vector q′+0 ,
which must be applied to the interceptor just after the first firing of the rocket motor
in order for the interceptor to reach the point of rendezvous after the lapse of a
specified time tf from the initial time t = 0.
Of course, the components x′−0 , y′−0 , and z′−0 of the velocity vector q′−0 applied to
the interceptor just before the first firing are known, because the orbital elements of
the parking orbit of the interceptor are known.
The first variation of velocity, which is needed by the interceptor at the initial
time t = 0, results from the following vector expression
þ
Dq00 ¼ q0 0 q0
0
At the time tf, the interceptor must arrive at the point of rendezvous, that is, at the
origin O of the non-inertial reference system Oxyz where the centre of mass of the
target is located. In other words, the position vector qf of the interceptor at the time
tf, with respect to the non-inertial reference system Oxyz, must be equal to 0.
This condition is expressible in mathematical terms as follows
qf qðtf Þ ¼ 0
On the other hand, by using the analytical solution of the Hill–Clohessy–

Wiltshire equations derived in Sect. 4.10, the position vector q(t) of the interceptor
at any time t, with respect to Oxyz, may be written as follows
Therefore, the preceding condition becomes

þ
qf qðtf Þ ¼ Urr ðtf Þq0 þ Urv ðtf Þq0 0 ¼0
The preceding equation, solved for q′+0 , yields
q00þ ¼ ½Urv ðtf Þ1 ½Urr ðtf Þq0
where the matrix [Urv(tf)]−1 is the inverse of Urv(tf). The preceding equation makes
it possible to compute the velocity vector q′+0 of the interceptor at the beginning of
the transfer ellipse. By substituting the preceding equation
q00þ ¼ ½Urv ðtf Þ1 ½Urr ðtf Þq0
into the solution (see Sect. 4.10) of the Hill–Clohessy–Wiltshire equations
q0 ðtÞ ¼ Uvr ðtÞq0 þ Uvv ðtÞq00

with t = t−f , it is possible to determine the velocity vector q′−f which must be applied
to the interceptor at the point of its arrival at the target, as follows
þ
q0 0
f q tf ¼ Uvr tf q0 þ Uvv tf q0 0
n 1 o
¼ Uvr tf q0 þ Uvv tf Urv tf Urr tf q0
By simplifying the preceding equation, there results

n 1 o
q0
f ¼ U t
vr f

U vv tf Urv tf Urr tf q0
When the interceptor has reached the target, it is necessary to perform the second
variation of velocity, in order for the two vehicles to have the same velocity after
the second burn, that is, at time t = t+f . In mathematical terms, in order for the
velocity vector q′+f q′(t+f ) to be equal to 0, the second variation of velocity Dq′f
must be such that
Dq0f q0f þ q0 0 0

f ¼ 0 qf ¼ qf
The following example will show the sequence of calculations to be performed

in practice. At a given epoch t = 0, the classical orbital elements of a space station
are the following
a ¼ 6788:0 km e ¼ 0:0 i ¼ 51
:647
X ¼ 278
:783 x ¼ 0 ðarbitraryÞ M0 ¼ 326
:384
The corresponding position and velocity vectors of the space station with respect
to the geocentric equatorial reference system XYZ can be computed as has been
shown in Sect. 1.9. This computation gives the following results:
r0 ¼ 1441:4 uX 5942:6 uY 2947:0 uZ ðkmÞ

v0 ¼ 4:5611 uX 3:5880 uY þ 5:0044 uZ ðkm=sÞ
At the same epoch, the classical orbital elements of an interceptor are the
following
a ¼ 7263:4 km e ¼ 0:065211 i ¼ 50
:845
X ¼ 99
:498 x ¼ 155
:79 M0 ¼ 9
:8197
The corresponding position and velocity vectors of the interceptor with respect
to the geocentric equatorial reference system are:
r ¼ 1440:3 uX 5950:2 uY 2950:4 uZ ðkmÞ

v ¼ 4:8243 uX 3:6082 uY þ 5:1123 uZ ðkm=sÞ
Let Oxyz be the radial-transverse-bi-normal system of reference, whose origin O

coincides with the instantaneous position of the centre of mass of the space station.
According to the definition given in Sect. 4.10, the unit vectors of Oxyz are such
that
r0 v0
ux ¼ uy ¼ uz ¼ ux uy
jr0 j j v0 j
Hence
h i12
jr0 j ¼ ð1441:4Þ2 þ ð5942:6Þ2 þ ð2947:0Þ2 ¼ 6788:0 km
h i12
jv0 j ¼ 4:56112 þ ð3:5880Þ2 þ 5:00442 ¼ 7:6630 km/s
1441:4 5942:6 2947:0

ux ¼ uX uY uZ
6788:0 6788:0 6788:0
¼ 0:21235uX 0:87546uY 0:43415uZ
4:5611 3:5880 5:0044
uy ¼ uX uY þ uZ
7:6630 7:6630 7:6630
¼ 0:59521uX 0:46822uY þ 0:65306uZ
uz ¼ ð0:87546 0:65306 0:46822 0:43415ÞuX þ ð0:21235 0:65306
0:59521 0:43415ÞuY þ ð0:21235 0:0:46822 þ 0:59521 0:87546ÞuZ
¼ 0:77501uX 0:11973uY þ 0:62051uZ
Therefore, the rotation matrix, to be used to transform the XYZ co-ordinates into
xyz co-ordinates, is
2 3
0:21235 0:87546 0:43415
R 4 0:59521 0:46822 0:65306 5
0:77501 0:11973 0:62051
On the other hand, in the geocentric equatorial system, the position vector q of
the interceptor relative to the space station is
q ¼ r r0 ¼ ð1440:3 þ 1441:4Þ uX þ ð5950:2 þ 5942:6Þ uY

þ ð2950:4 þ 2947:0Þ uZ ¼ 1:1 uX 7:6 uY 3:4 uZ
In the geocentric equatorial system, the velocity vector q′ of the interceptor

relative to the space station is
q0 ¼ v v0 x q
where x = x uz, and x is the mean motion of the space station, resulting from
!12 1
lE 398600:44 2
x¼ ¼ ¼ 0:0011289 rad/s
jr0 j3 6766:03
Hence
x ¼ 0:0011289ð0:77501 uX 0:11973 uY þ 0:62051 uZ Þ

¼ 0:00087491 uX 0:00013516 uY þ 0:00070049 uZ
x q ¼ ð0:00013516 3:4 þ 0:70049 7:6Þ uX ð0:00087491 3:4 0:00070049

1:1Þ uY þ ð0:00087491 7:6 þ 0:00013616 1:1Þ uZ
¼ 0:0057833 uX 0:0022042 uY þ 0:0067980 uZ
q0 ¼ v v0 x q ¼ ð4:8243 4:5611 0:0057833Þ uX þ ð3:6082 þ 3:5880

þ 0:0022042Þ uY þ ð5:1123 5:0044 0:0067980Þ uZ
¼ 0:25742 uX 0:017996 uY þ 0:10110 uZ
Now, the vectors q and q′ computed above in the geocentric equatorial system
XYZ are expressed in the system of reference Oxyz of the space station by using the
rotation matrix R. This yields
2 3 2 32 3 2 3
x 0:21235 0:87546 0:43415 1:1 7:8960
6 7 6 76 7 6 7
¼
4 5 4 0:59521
y 0:46822 0:65306 54 7:6 5 ¼ 4 1:9928 5
z 0:77501 0:11973 0:62051 3:4 2:0523
2 03 2 32 3 2 3
x 0:21235 0:87546 0:43415 0:25742 0:082801
6 07 6 76 7 6 7
4 y 5 ¼ 4 0:59521 0:46822 0:65306 54 0:017996 5 ¼ 4 0:22767 5
z0 0:77501 0:11973 0:62051 0:10110 0:13461
The co-ordinates x, y, and z are the components of the position vector q0 of the
interceptor just before the first firing of the rocket motor. Likewise, x′, y′, and z′ are
the components of the velocity vector q′−0 just before the first firing of the rocket
motor.
Now, the matrices Urr(t), Urv(t), Uvr(t), and Uvv(t) defined in Sect. 4.10 are
computed for x = 0.0011289 rad/s and a final time t = tf, which we choose equal to
1 h = 3600 s. This yields
2 3 2 3
4 3 cos x t 0 0 5:8116 0 0
6 7 6 7
Urr ðtÞ 4 6ðsin x t x tÞ 1 0 ¼
5 4 29:167 1 0 5
0 0 cos x t 0 0 0:60387
2 3
ð1=xÞ sin x t ð2=xÞð1 cos x tÞ 0
6 7
Urv ðtÞ 4 ð2=xÞðcos x t 1Þ ð4=xÞ sin x t 3t 0 5
0 0 ð1=xÞ sin x t
2 3
706:07 2841:5 0
6 7
¼ 4 2841:5 13624 0 5
0 0 706:07
2 3 2 3
3x sin x t 0 0 0:0026995 0 0
6 7 6 7
Uvr ðtÞ 4 6xðcos x t 1Þ 0 0 5 ¼ 4 0:010864 0 0 5
0 0 x sin x t 0 0 0:00089983
2 3
cos x t 2 sin x t 0
6 7
Uvv ðtÞ 4 2 sin x t 4 cos x t 3 0 5
0 0 cos x t
2 3
0:60387 1:5942 0
6 7
¼ 4 1:5942 5:4155 0 5
0 0 0:60387
The velocity q′+0 of the interceptor at the beginning of the transfer ellipse is
þ
q0 0 ¼ ½Urv ðtf Þ1 ½Urr ðtf Þq0

2 3
0:00077000 0:00016060 0
6 7
½Urv ðtf Þ1 ¼ 4 0:00016060 0:000039905 0 5
0 0 0:0014163
2 32 3 2 3
5:8116 0 0 7:8960 45:888
6 76 7 6 7
Urr ðtf Þq0 ¼ 4 29:167 1 0 54 1:9928 5 ¼ 4 228:31 5
0 0 0:60387 2:0523 1:2393
2 32 3 2 3
0:00077000 0:00016060 0 45:888 0:0013328
þ 6 76 7 6 7
q0 0 ¼ 4 0:00016060 0:000039905 0 54 228:31 5 ¼ 4 0:016480 5
0 0 0:0014163 1:2393 0:0015552
The velocity vector q′−f which must be applied to the interceptor at the point of
its arrival at the target results from
n 1 o
q0
f ¼ Uvr tf Uvv tf Urv tf Urr tf q0

2 3
0:0026995 0 0
6 7
Uvr tf 4 0:010864 0 0 5
0 0 0:00089983
2 3
0:00077000 0:00016060 0
1 6 7
½Urv ðtf Þ 4 0:00016060 0:000039905 0 5
0 0 0:0014163
2 3
5:8116 0 0
6 7
Urr tf 4 29:167 1 0 5
0 0 0:60387
2 3
0:00020929 0:00016060 0
1 6 7
Urv tf Urr tf 4 0:0020973 0:000039905 0 5
0 0 0:00085526
2 3
0:60387 1:5942 0
6 7
Uvv ðtf Þ 4 1:5942 5:4155 0 5
0 0 0:60387
2 3
0:0046990 0:00016060 0
1 6 7
Uvv tf Urv tf Urr tf ¼ 4 0:011024 0:000039923 0 5
0 0 0:00051647
2 3
0:0020291 0:0001606 0
1 6 7
Uvr tf Uvv tf Urv tf Urr tf ¼ 4 0:00016 0:000039923 0 5
0 0 0:0014163
2 32 3
0:0020291 0:0001606 0 7:8960
6 76 7
q0
f ¼ 4 0:00016 0:000039923 0 54 1:9928 5
0 0 0:0014163 2:0523
2 3
0:015702
6 7
¼ 4 0:0013429 5
0:0029067
The first variation of velocity, which is needed by the interceptor at the initial
time t = 0, results from
þ
Dq00 ¼ q0 0 q0
0
By introducing the values computed above for q′+0 and q′−0 , there results
2 3 2 3 2 3
0:0013328 0:015702 0:15835
Dq00 ¼ q00þ q0
0 ¼ 4 0:016480 5 4 0:0013429 5 ¼ 4 0:017823 5
0:0015552 0:0029067 0:0044619
The second variation of velocity Dq′f, to be applied when the interceptor has
reached the target, results from
Dq0f q0f þ q0 0 0

f ¼ 0 qf ¼ qf
By introducing the values computed above for q′−f , there results

2 3 2 3 2 3
0 0:015702 0:015702
Dq0f ¼ 0 q0
f ¼ 4 0 5 4 0:0013429 5 ¼ 4 0:0013429 5
0 0:0029067 0:0029067
The magnitudes of these variations of velocity are respectively
0 h i1
Dq ¼ ð0:15835Þ2 þ ð0:017823Þ2 þ 0:00446192 2 ¼ 0:15941 km/s
0
0 h i1
Dq ¼ ð0:015702Þ2 þ ð0:0013429Þ2 þ 0:00290672 2 ¼ 0:016025 km/s
f
The magnitude of the total variation of velocity necessary to perform the ren-
dezvous manoeuvre results from
0 0
Dq þ Dq ¼ 0:15941 þ 0:016025 ¼ 0:17544 km=s ¼ 175:44 m=s
0 f
Within the time interval 0 < t < tf, the position vector q(t) with respect to Oxyz
results from the following equation derived in Sect. 4.10:
where the matrices Urr(t) and Urv(t) are in the general case
2 3
4 3 cos x t 0 0
Urr ðtÞ 4 6ðsin x t x tÞ 1 0 5
0 0 cos x t
2 3
ð1=xÞ sin x t ð2=xÞð1 cos x tÞ 0
Urv ðtÞ 4 ð2=xÞðcos x t 1Þ ð4=xÞ sin x t 3t 0 5
0 0 ð1=xÞ sin x t
and, in the present example, the position and velocity vectors q0 and q′0 have the
following components
2 3 2 3
7:8960 0:082801
q0 4 1:9928 5 q00 4 0:22767 5
2:0523 0:13461
By substituting 0.0011289 in place of x in the matrices Urr(t) and Urv(t), it is

possible to compute the relative position vector q(t) as a function of time.
4.12 Hohmann Transfer Manoeuvres
The fundamental concepts on these manoeuvres have been given in Sect. 4.6. The
present paragraph and those which follow will show the matter in detail. Practical
examples of application will also be given.
A Hohmann transfer manoeuvre [53] consists of two impulses, the first of which
produces a transfer ellipse between two circular and coplanar orbits of different radii
around a common centre of attraction, and the second produces a circular orbit of
higher radius by using the least amount of propulsive energy.
It is an ideal transfer manoeuvre, which cannot be performed in practice, because
of unavoidable errors in the magnitudes and in the times of application of the
impulses which are actually imparted to a spacecraft by means of its rocket motor.
4.12 Hohmann Transfer Manoeuvres 685
With reference to the preceding figure, taken from Wikimedia [71], it is desired
to move a spacecraft from a circular orbit (1) of radius R r1 around the Earth to
another circular orbit (3) of radius R′ r2, such that r2 > r1, by means of an elliptic
transfer orbit (2) which touches the low circular orbit (1) at perigee and the high
circular orbit (3) at apogee.
Since r1 is the radius at perigee and r2 is the radius at apogee, then the major
axis, 2a, of the elliptic transfer orbit must be such that
2a ¼ r1 þ r2
For this purpose, the rocket motor of the spacecraft is first fired in such a way as
to provide a velocity increment Dv Dv1, which brings the spacecraft from the low
circular orbit (1) to the elliptic transfer orbit (2), and then fired again at the apogee
of the transfer orbit (2), by means of a velocity increment Dv′ Dv2, in order for
the spacecraft to move from the elliptic transfer orbit (2) to the circular orbit (3).
A Hohmann transfer manoeuvre may also be used to bring a spacecraft from a
high circular orbit of radius r1 to a low circular orbit of radius r2 (with r2 < r1), but
in this case the rocket motor must be fired in the direction opposite to that of the
flight path, in order for the spacecraft to approach the centre of attraction. Such is
the case with an interplanetary probe which is to be brought from the orbit of the
Earth to the orbit of Venus, as shown in the following figure, due to the courtesy of
NASA-JPL [72].
The total variation of velocity needed for the manoeuvre results from the sum of
the two variations Dv1 and Dv2. Of course, the expenditure of propellant depends on
the sum of the absolute values of Dv1 and Dv2 in both cases, independently of the
direction along which the rocket motor is fired.
Under the hypotheses indicated above (circular orbits of departure and arrival,
elliptic transfer orbit going from perigee to apogee, and impulsive firings), we want
to compute the two variations of velocity, Dv1 and Dv2, to be performed at,
respectively, perigee and apogee of the elliptic transfer orbit.
To this end, we use the well-known energy integral

2 1
v2 ¼ lE
r a
where v and r are, respectively, the magnitude of the velocity and the radius vector
of the spacecraft at some point of its transfer ellipse, lE is the gravitational
parameter of the Earth, and a = ½(r1 + r2) is the major semi-axis of the transfer
ellipse from the low-altitude circular orbit of radius r1 to the high-altitude circular
orbit of radius r2.
Let v−1 and v+1 be the magnitudes of the velocities, respectively, immediately
before and immediately after the time in which the first impulse Dv1 is applied to
the spacecraft, such that Dv1 = v+1 − v−1 . Since Dv1 is applied at the perigee (r1) of
the transfer ellipse, there results
12
lE
v
1 ¼
r1
12
2lE 2lE
v1þ ¼
r1 r1 þ r2
By executing the subtraction v+1 − v−1 and simplifying, there results

12 " 12 #
lE 2r2
Dv1 ¼ 1
r1 r1 þ r2
Likewise, let v−2 and v+2 be the magnitudes of the velocities, respectively,
immediately before and immediately after the time in which the second impulse Dv2
is applied to the spacecraft, such that Dv2 = v+2 − v−2 . Since Dv2 is applied at the
apogee (r2) of the transfer ellipse, there results
1
2lE 2lE 2
v
2 ¼
r2 r1 þ r2
12
lE
v2þ ¼
r2
By executing the subtraction v+2 − v−2 and simplifying, there results

12 " 1 #
lE 2r1 2
Dv2 ¼ 1
r2 r1 þ r2
The total variation of velocity needed for a Hohmann transfer manoeuvre results
from the sum of the two variations Dv1 and Dv2 computed above.
The time, t2 − t1, needed to perform this manoeuvre is equal to the semi-period
of the transfer ellipse, that is,
" 1 # 3 12
1 1 a3 2 a
t2 t1 ¼ T ¼ 2p ¼p
2 2 lE lE
This is because a spacecraft travels through an angle of p radians when a

Hohmann transfer manoeuvre is performed. Since the major semi-axis of the
transfer ellipse is a = ½(r1 + r2), then
" #12
1 ðr1 þ r2 Þ3
t2 t1 ¼ p
8 lE
Of course, the equation given above holds independently of whether a spacecraft

goes from a low-altitude circular orbit to a high-altitude circular orbit or vice versa.
When a Hohmann transfer manoeuvre is performed in a travel from one celestial
body (e.g. from the Earth) to another (e.g. to Venus), it is of paramount importance
to start the manoeuvre when the lead angle w (in radians) between the two planets at
the moment of the spacecraft insertion into the transfer semi-ellipse (which angle is
illustrated in the first figure of Sect. 4.9) has the proper value. In order for the
spacecraft and the target planet to be in the same position at the final time of a
Hohmann transfer manoeuvre, the following condition must be satisfied
w þ nP ðt2 t1 Þ ¼ p
where nP = (lS/a3P)½ is the mean motion of the target planet, lS is the gravitational
parameter of the Sun, aP is the major semi-axis of the orbit of the target planet, and
p radians is the travel angle of the spacecraft between the orbits of the two planets.
By substituting the expression of t2 − t1 computed above into the preceding
equation and solving for w, with aP = r2, there results
8
" 9
< 3 #12 =
1 r1
w¼p 1 þ1
: 8 r2 ;
This determines the launch time, because an Earth-to-Venus travel based on a

Hohmann transfer manoeuvre is only possible when the lead angle w has the value
given above, which fact occurs every 19 months [72]. In case of an Earth-to-Mars
travel, the same opportunity occurs every 25 months [72].
By the way, the first of the three methods described in Sect. 4.9 does not require
a travel angle of p radians between the two circular orbits of departure and arrival,
that is, this method can be used independently of whether the transfer manoeuvre
performed by a spacecraft be, or be not, a Hohmann transfer manoeuvre. The latter
has over the former the advantage of requiring a lower amount of propulsive
energy. We show this fact by means of the following example. Let us consider a
spacecraft revolving about the Earth along a circular equatorial orbit of radius
r1 = 6678 km, corresponding to an altitude h1 = 300 km. We want to raise this
spacecraft to a geostationary orbit of radius r2 = 42164 km. To this end, we use two
different methods: (1) a Hohmann transfer manoeuvre, and (2) a higher-energy
transfer manoeuvre, in order to obtain the same result in less time, by spending a
greater amount of propulsive energy.
Method (1): Hohmann transfer manoeuvre.
The major semi-axis of the transfer semi-ellipse results from
r1 þ r2 6678 þ 42164
aH ¼ ¼ ¼ 24421 km
2 2
where the subscript H indicates a Hohmann transfer manoeuvre.

The first increment of velocity, needed to inject the spacecraft from the circular
orbit of radius r1 into the transfer semi-ellipse of major semi-axis aH, results from
12 " 12 #
lE 2r2
Dv1 ¼ 1
r1 r1 þ r2

1 " 1 #
398600 2 2 42164 2
Dv1 ¼ 1 ¼ 2:4258 km/s
6678 6678 þ 42164
The second increment of velocity, needed to inject the spacecraft from the
transfer semi-ellipse of major semi-axis aH into the circular orbit of radius r2, results
from
12 " 1 #
lE 2r1 2
Dv2 ¼ 1
r2 r1 þ r2

1 " 12 #
398600 2 2 6678
Dv2 ¼ 1 ¼ 1:4668 km/s
42164 6678 þ 42164
The total variation of velocity results from the sum of the two variations Dv1 and
Dv2 computed above. In the present case, we have
Dv ¼ Dv1 þ Dv2 ¼ 2:4258 þ 1:4668 ¼ 3:8926 km=s
The time needed to perform a Hohmann transfer manoeuvre is

" #12
1 ðr1 þ r2 Þ3
t2 t1 ¼ p
8 lE

" #12
ð6678 þ 42164Þ3
t2 t1 ¼ 3:1416 ¼ 18990s ¼ 5:2750 h
8 398600
Method (2): High-energy transfer orbit.

We choose arbitrarily the major semi-axis of the transfer trajectory as follows
aTR ¼ 24500 km
which value is greater than aH = 24421 km. Of course, if the value of aTR were
chosen less than 24421 km, the transfer trajectory would not intersect the circular
orbit of radius r2. The choice made above assures us that the transfer trajectory
intersects this circular orbit in two distinct points. Let R and S be such points. We
search the point R of intersection whose true anomaly, /R, is less than p radians.
The subscript TR indicates here the transfer trajectory between the two circular
orbits. The eccentricity of the transfer trajectory results from
r1 6678
eTR ¼ 1 ¼1 ¼ 0:72743
aTR 24500
The semi-latus rectum of the transfer trajectory results from

pTR ¼ aTR 1 e2TR ¼ 24500 1 0:727432 ¼ 11536 km
The spacecraft is inserted into the transfer trajectory at perigee (/ = 0). At the
point R of intersection between the transfer trajectory and the circular orbit of radius
r2 = 42164 km, the true anomaly, that is, the travel angle of the spacecraft, com-
puted from the perigee of the transfer trajectory to the point R of intersection
between the transfer trajectory and the circular orbit of radius r2, is (in radians):
0p 1 0 1
TR 11536
1 1
Br C B C
/R ¼ arccos@ 2 A ¼ arccos@42164 A ¼ 3:0886
eTR 0:72743
which value is less than p radians.

Since the flight path angle, c, at any point of an elliptic orbit is by definition
vr e sin /
tan c ¼ ¼
v/ 1 þ e cos /
then the flight path angle (in radians) at the point R of intersection between the
transfer trajectory and the circular orbit of radius r2 = 42164 km is

e sin / 0:72743 sin 3:0886
cR ¼ arctan ¼ arctan ¼ 0:13993
1 þ e cos / 1 þ 0:72743 cos 3:0886
The difference of velocity (Dv1), which is needed at the moment in which the
spacecraft is inserted from the circular orbit of radius r1 = 6678 km into the transfer
trajectory, is
12 12
2 1 l
Dv1 ¼ lE E
r1 aTR r1

12 1
2 1 398600 2
Dv1 ¼ 398600 ¼ 2:4284 km/s
6678 24500 6678
At the point R of intersection between the transfer trajectory and the circular
orbit of radius r2 = 42164 km, the magnitude, vR, of the velocity vector vR tangent
to the transfer trajectory is
12
2 1
vR ¼ lE
r2 aTR

12
2 1
vR ¼ 398600 ¼ 1:6241 km/s
42164 24500
The magnitude, vT, of the velocity vector vT tangent to the circular orbit of radius
r2 = 42164 km is
12
lE
vT ¼
r2

12
398600
vT ¼ ¼ 3:0747 km/s
42164
Since cR = 0.13993 rad is the angle between the two velocity vectors vR and vT,
then the difference of velocity (Dv2), which is needed at the moment in which the
spacecraft is inserted from the transfer trajectory into the circular orbit of radius
r2 = 42164 km, results from the cosine rule
1
Dv2 ¼ v2R þ v2T 2vR vT cos cR 2

1
Dv2 ¼ 1:62412 þ 3:07472 2 1:6241 3:0747 cos 0:13993 2 ¼ 1:4838 km/s
The total difference of velocity is therefore
Dv ¼ Dv1 þ Dv2 ¼ 2:4284 þ 1:4838 ¼ 3:9122 km=s
which value is higher than the value of 3.8926 km/s computed above for a
Hohmann transfer manoeuvre.
By using the following expression of Sect. 1.3:
where e = eTR and / = /R, we compute the eccentric anomaly, in radians, at the
point R of intersection, as follows
Likewise, by using the following expression of Sect. 1.3:
where e = eTR and Æ = ÆR, it is possible to compute the mean anomaly, in radians,
at the point R of intersection as follows
MR ¼ 3:0083 0:72743 sin 3:0083 ¼ 2:9117
The transfer time of the spacecraft along the transfer trajectory, from the point of
perigee to the point R of intersection with the circular orbit of radius r2, results from
3 12
aTR
t2 t1 ¼ MR
lE

12
245003
t2 t1 ¼ 2:9117 ¼ 17686 s ¼ 4:9127 h
398600
which value is lower than the value of 5.2750 h computed above for a Hohmann
transfer manoeuvre.
By the way, the transfer time t2 − t1 may also be computed without using the
eccentric anomaly corresponding to the true anomaly /R, by first evaluating the
chord c joining the starting point (perigee) and the intersection point R, as shown
below
1
c ¼ r12 þ r22 2r1 r2 cos /R 2
Then, by applying Lambert’s theorem (see Sect. 1.5), the transfer time results
from
3 12
aTR
t2 t1 ¼ ½ða sin aÞ ðb sin bÞ
lE
where a and b are, respectively, Lambert’s first and second angle, expressed by
a r þ r þ c12
b
1
r1 þ r2 c 2
1 2
sin ¼ sin ¼
2 4aTR 2 4aTR
By so doing, we have
1
c ¼ 66782 þ 421642 2 6678 42164 cos 3:0886 2 ¼ 48834 km
" 12 #
6678 þ 42164 þ 48834
a ¼ 2 arcsin ¼ 3:0265 rad
4 24500
" 12 #
6678 þ 42164 48834
b ¼ 2 arcsin ¼ 0:018178 rad
4 24500
12
245003
t2 t1 ¼ ½ð3:0265 sin 3:0265Þ ð0:018178 sin 0:018178Þ
398600
¼ 17686 s ¼ 4:9127 h
As was expected, a Hohmann transfer manoeuvre (which uses a transfer orbit of

major axis 2aH = r1 + r2) requires a smaller amount of propulsive energy than that
required by another transfer manoeuvre, which uses a transfer orbit of major axis
2aTR > r1 + r2. These manoeuvres are meant both of them to bring a spacecraft
from a circular orbit of radius r1 to another circular orbit of radius r2.
So far, a Hohmann transfer manoeuvre has been defined with reference to a
circular orbit of departure and a circular orbit of arrival. However, a Hohmann
transfer manoeuvre can also be performed when the two orbits of departure and of
arrival are elliptic instead of circular, provided that the two ellipses have the same
apsidal line. The example given below will show the practical procedure.
With reference to the following figure (where the image of the Earth is due to the
courtesy of NASA), we consider a spacecraft revolving about the Earth in an
equatorial (i = 0) orbit of low altitude. Let r1 and r2 be, respectively, the radius of
perigee and the radius of apogee of the low-altitude orbit measured from the centre
of mass of the Earth. We want to use a Hohmann transfer manoeuvre to move the
spacecraft from the low-altitude elliptic orbit described above to a high-altitude
elliptic orbit, whose apsidal line coincides with the apsidal line of the low-altitude
orbit.
Let R1 and R2 be, respectively, the radius of perigee and the radius of apogee of
the high-altitude orbit measured from the centre of mass of the Earth. The direction
of motion for both orbits is counterclockwise, for an observer placed above the
north pole of the Earth.
In the present example, the Hohmann transfer manoeuvre will be considered in
two different cases:
(1) an impulse given to the spacecraft at the apogee of the low-altitude ellipse, such
as to insert the spacecraft into a transfer semi-ellipse ending at the perigee of
this semi-ellipse (which point is also the perigee of the high-altitude ellipse),
followed by another impulse, given to the spacecraft at this point, such as to
insert the spacecraft into the final high-altitude ellipse; or
(2) an impulse given to the spacecraft at the perigee of the low-altitude ellipse, such
as to insert the spacecraft into a transfer semi-ellipse ending at the apogee of
this semi-ellipse (which point is also the apogee of the high-altitude ellipse),
followed by another impulse, given to the spacecraft at this point, such as to
insert the spacecraft into the final high-altitude ellipse.
The choice between a Hohmann transfer of the first type and a Hohmann transfer
of the second type depends on which of the two types of manoeuvre requires less
propulsive energy than the other type, as the sequel will show.
We give the following values to the radii of the two ellipses indicated above:
r1 ¼ 6644 km r2 ¼ 11428 km
R1 ¼ 6910 km R2 ¼ 19309 km
Hence, the major semi-axes of the two elliptic orbits are
r1 þ r2 6644 þ 11428
a¼ ¼ ¼ 9036 km
2 2
R1 þ R2 6910 þ 19309
A¼ ¼ ¼ 13109:5 km
2 2
(1) Hohmann transfer of the first type

The velocity possessed by the spacecraft at the apogee of the low-altitude orbit,
just before the first impulse, is
1 12
2 1 2 2 1
v
1 ¼ lE ¼ 398600 ¼ 5:0642 km/s
r2 a 11428 9036
At the same point, the velocity, which the spacecraft must possess in order to be
inserted into a transfer semi-ellipse of major semi-axis
aH = ½(r2 + R1) = 9169 km, just after the first impulse, is
12 12
2 1 2 1
v1þ ¼ lE ¼ 398600 ¼ 5:1270 km/s
r 2 aH 11428 9169
The first increment of velocity to be given to the spacecraft is then
Dv1 ¼ v1þ v
1 ¼ 5:1270 5:0642 ¼ 0:06280 km=s
The velocity possessed by the spacecraft at the perigee of the transfer

semi-ellipse of major semi-axis aH, just before the second impulse, is
12 12
2 1 2 1
v
2 ¼ lE ¼ 398600 ¼ 8:4792 km/s
R 1 aH 6910 9169
inserted into the high-altitude ellipse of major semi-axis A = 13109.5 km, just after
the second impulse, is
1 12
2 1 2 2 1
v2þ ¼ lE ¼ 398600 ¼ 9:2176 km/s
R1 A 6910 13109:5
The second increment of velocity to be given to the spacecraft is then
Dv2 ¼ v2þ v
2 ¼ 9:2176 8:4792 ¼ 0:7384 km=s
The total increment of velocity for a Hohmann transfer of the first type is then
Dv ¼ Dv1 þ Dv2 ¼ 0:06280 þ 0:7384 ¼ 0:8012 km=s
(2) Hohmann transfer of the second type

The velocity possessed by the spacecraft at the perigee of the low-altitude orbit,
just before the first impulse, is
12 12
2 1 2 1
v
1 ¼ lE ¼ 398600 ¼ 8:7107 km/s
r1 a 6644 9036
At the same point, the velocity, which the spacecraft must possess in order
to be inserted into a transfer semi-ellipse of semi major axis
aH = ½(r1 + R2) = 12976.5 km, just after the first impulse, is
12 12
2 1 2 1
v1þ ¼ lE ¼ 398600 ¼ 9:4483 km/s
r 1 aH 6644 12976:5
The first increment of velocity to be given to the spacecraft is then
Dv1 ¼ v1þ v
1 ¼ 9:4483 8:7107 ¼ 0:7376 km=s
The velocity possessed by the spacecraft at the apogee of the transfer

semi-ellipse of major semi-axis aH, just before the second impulse, is
12 12
2 1 2 1
v
2 ¼ lE ¼ 398600 ¼ 3:2511 km/s
R2 aH 19309 12976:5
inserted into the high-altitude ellipse of major semi-axis A = 13109.5 km, just after
the second impulse, is
12 12
2 1 2 1
v2þ ¼ lE ¼ 398600 ¼ 3:2986 km/s
R2 A 19309 13109:5
The second increment of velocity to be given to the spacecraft is then
Dv2 ¼ v2þ v
2 ¼ 3:2986 3:2511 ¼ 0:0475 km=s
The total increment of velocity for a Hohmann transfer of the second type is then
Dv ¼ Dv1 þ Dv2 ¼ 0:7376 þ 0:0475 ¼ 0:7851 km=s
This shows that a Hohmann transfer manoeuvre of the second type requires less
propulsive energy than one of the first type, because 0.7851 < 0.8012. This is
because the velocity of an orbiting body is maximum at perigee and minimum at
apogee, according to the second law of Kepler. This is also in compliance with the
Oberth effect (so called after the German scientist Hermann Oberth): a rocket motor
generates, when travelling at high velocity, a greater amount of useful energy than
that generated at low velocity.
4.13 Bi-Elliptic Transfer Manoeuvres
A bi-elliptic transfer manoeuvre is a sequence of two Hohmann transfer manoeu-

vres, which have been described in Sect. 4.12. Under certain conditions, which will
be specified below, a bi-elliptic transfer manoeuvre requires a smaller amount of
propulsive energy than that required by a single Hohmann transfer manoeuvre.
On the other hand, the time necessary for a bi-elliptic transfer manoeuvre is
considerably greater than the time for a single Hohmann transfer manoeuvre.
A bi-elliptic transfer manoeuvre is based on the method of moving a spacecraft so
far away from the centre of mass of the main attracting body as to obtain the desired
effect by means a small variation of velocity. With reference to the following figure,
taken from Wikimedia [73], a bi-elliptic transfer manoeuvre starts from a
low-altitude circular orbit (blue) around the Sun and ends to a high-altitude circular
orbit (red) around the Sun.
For this purpose, the spacecraft is moved, by means of a variation of velocity

Dv1, at a high distance rb from the centre of mass of the Sun. In the following
section of the present paragraph, we use the notation of Sect. 4.12, that is, r1 is the
radius of the circular orbit (blue) of departure and r2 is the radius of the circular
orbit (red) of arrival. Three variations of velocity are necessary for this manoeuvre.
The first (Dv1) inserts the spacecraft from the low-altitude circular orbit (blue) of
departure into a transfer semi-ellipse (cyan) of major semi-axis aH1 = ½(r1 + rb);
the second (Dv2) inserts the spacecraft from this transfer semi-ellipse into another
transfer semi-ellipse (orange) of major semi-axis aH2 = ½(r2 + rb); and the third
(Dv3), opposite to the orbital velocity, inserts the spacecraft from the second transfer
semi-ellipse (orange) into the high-altitude circular orbit (red) of radius r2.
Let the Earth be the main attracting body. By using repeatedly the integral of
energy v2 = lE(2/r − 1/a), it is easy to compute each of the three increments of
velocity, as follows
12 12
2 1 l
Dv1 ¼ lE E
r1 aH1 r1
12 12
2 1 2 1
Dv2 ¼ lE lE
rb aH2 rb aH1
12 12
2 1 l
Dv3 ¼ lE E
r2 aH2 r2
The total variation of velocity required for a bi-elliptic transfer manoeuvre is
ðDvÞbielliptic ¼ Dv1 þ Dv2 þ Dv3
which is the sum of the three increments indicated above, independently of the
direction of each single variation of velocity.
The time needed to perform a bi-elliptic transfer manoeuvre is equal to the sum
of the two semi-periods of the two transfer ellipses, that is,
3 12 3 12
aH1 a
ðt2 t1 Þbielliptic ¼ p þ p H2
lE lE
where, in all of the equations written above, aH1 = ½(r1 + rb) and
aH2 = ½(r2 + rb).
On the other hand, as has been shown in Sect. 4.12, the two velocity increments
Dv1 and Dv2 relating to a Hohmann transfer manoeuvre between the same initial
and final circular orbits are
12 " 1 #
lE 2r2 2
Dv1 ¼ 1
r1 r1 þ r2
12 " 1 #
lE 2r1 2
Dv2 ¼ 1
r2 r1 þ r2
4.13 Bi-Elliptic Transfer Manoeuvres 699
The total variation of velocity required for a Hohmann transfer manoeuvre is
ðDvÞHohmann ¼ Dv1 þ Dv2
which is the sum of the two velocity increments Dv1 and Dv2 indicated above.
A bi-elliptic transfer manoeuvre is advantageous over a Hohmann transfer
manoeuvre having the same circular orbits of, respectively, departure and arrival
when the total difference of velocity required by the former is lower than that
required by the latter, that is, when
ðDvÞbielliptic \ðDvÞHohmann
Escobal [74] has shown that the preceding condition

• is satisfied when r2 is greater than about 15.58 times r1;
• is not satisfied when r2 is less than about 11.94 times r1; and
• is, or is not, satisfied when r2 falls between 11.94 r1 and 15.58 r1, depending on
whether rb is, or is not, much greater than r1.
For example, if the radii of the two circular orbits were, respectively, r1 = 6678 km
and r2 = 42164 km, a bi-elliptic transfer manoeuvre would not be advantageous over a
Hohmann transfer manoeuvre, because 6678 11.94 = 79735.32, which value is
greater than 42164, let the value of rb be what it may.
The advantage of a bi-elliptic transfer manoeuvre over a Hohmann transfer
manoeuvre in terms of smaller propulsive energy is offset by the disadvantage of a
much longer time of flight, as the following example will show.
It is desired to move a spacecraft from a circular low orbit of radius
r1 = 6678 km around the Earth to a higher circular orbit of radius r2 = 86814 km.
Since r2 = 13r1, then it is possible to save a certain amount of propulsive energy by
using a bi-elliptic instead of a Hohmann transfer manoeuvre. The value of the
parameter rb must be chosen by the designer. We choose rb = 200340 km, that is,
30 times r1.
The major semi-axes of the two transfer semi-ellipses are
r1 þ rb 6678 þ 200340
aH1 ¼ ¼ ¼ 103509 km
2 2
r2 þ rb 86814 þ 200340
aH2 ¼ ¼ ¼ 143577 km
2 2
The three differences of velocity are

12 12 12 1
2 1 l 2 1 398600 2
Dv1 ¼ lE E ¼ 398600
r1 aH1 r1 6678 103509 6678
¼ 2:6101 km/s
12 12 12
2 1 2 1 2 1
Dv2 ¼ lE lE ¼ 398600
rb aH2 rb aH1 200340 143577
12
2 1
398600 ¼ 0:73855 km/s
200340 103509
2 2
1 1 12 1
2 1 lE 2 1 398600 2
Dv3 ¼ lE ¼ 398600
r2 aH2 r2 86814 143577 86814
¼ 0:38837 km/s
The total difference of velocity is
ðDvÞbielliptic ¼ Dv1 þ Dv2 þ Dv3 ¼ 2:6101 þ 0:73855 þ 0:38837 ¼ 3:7370 km=s
The time necessary to perform this manoeuvre is

3 12 3 12 1
aH1 aH2 1035093 2
ðt2 t1 Þbielliptic ¼p þp ¼ 3:1416
lE lE 398600
1
3 2
143577
þ 3:1416 ¼ 4:3642 105 s ¼ 121:23 h ¼ 5:0512 d
398600
By contrast, a Hohmann transfer manoeuvre leads to the following results (in

km/s)
12 " 12 # 1 " 1 #
lE 2r2 398600 2 2 86814 2
Dv1 ¼ 1 ¼ 1 ¼ 2:8029
r1 r1 þ r2 6678 6678 þ 86814
12 " 1 # 1 " 12 #

lE 2r1 2 398600 2 2 6678
Dv2 ¼ 1 ¼ 1
r2 r1 þ r2 86814 6678 þ 86814
¼ 2:1272
ðDvÞHohmann ¼ Dv1 þ Dv2 ¼ 2:8029 þ 2:1272 ¼ 4:0301 km=s

4.13 Bi-Elliptic Transfer Manoeuvres 701
The time necessary to perform this manoeuvre is

" #12 " #12
ðr1 þ r2 Þ3 ð6678 þ 86814Þ3
ðt2 t1 ÞHohmann ¼p ¼ 3:1416 ¼ 50292 s
8lE 8 398600
¼ 13:970 h ¼ 0:58208 d
The idea of a bi-elliptic transfer trajectory is due to Sternfeld [75]. The critical
coefficients 11.94 and 15.58 have been computed by Escobal [74].
4.14 Change of Orbital Plane
Manoeuvres to change the orbital plane of a satellite revolving about the Earth are
necessary when the latitude of the site of launch is greater than the desired angle of
the orbit with respect to the equatorial plane. For the reasons indicated above, it
may be necessary to change the inclination of the plane of motion of a satellite, in
such a way as to change the direction, but not the magnitude, of its velocity vector.
A simple plane change is a manoeuvre meant to that effect.
Firstly, as has been shown at length in Sects. 4.3 and 4.4, it is impossible to
insert directly a satellite into an orbit whose inclination, i, with respect to the
equator is less than the geocentric latitude, u*L, of the launch site. The equality
iMAX = u*L holds in case of a launch due East. Prograde (i < 90°) orbits having
inclinations 0 i < u*L require change-of-plane manoeuvres, to be performed
after the spacecraft has been inserted into a parking orbit of inclination i u*L. In
addition, a launch due East may be impossible owing to safety reasons.
Secondly, it has been shown that a rendezvous mission requires a specific inertial
orientation of the orbital plane, because, in order for the launched vehicle to reach
the orbital position of the target vehicle, the launch must occur around the time in
which the orbital plane of the target vehicle passes through the launch site. In case
of time discrepancies, a plane change is performed before a rendezvous manoeuvre,
in order for the orbital plane of the interceptor to be the same as the orbital plane of
the target.
Thirdly, the inertial velocity vector of spacecraft to be launched from a given site
may be subject to some requirements, in order for the spacecraft to accomplish its
mission, as is the case with lunar or planetary missions. As is well known, the
Moon and the other planets of the Solar System revolve in planes whose inclina-
tions are close to that of the ecliptic, but the geographic latitudes at which the
launch sites are located may be greater than these inclinations.
With reference to the preceding figure, let vi and vf be, respectively, the initial
velocity vector of a spacecraft before a simple plane change and the final velocity
vector of the same spacecraft after that manoeuvre. Let v be the intensity of both of
these vectors, such that |vi| = |vf| = v. Let Dv be the vector, whose intensity is Dv,
representing the increment of velocity to be applied to a spacecraft in order to
change the inclination of its orbital plane by an angle h.
For this purpose, a component of the velocity increment vector Dv must be
orthogonal to the initial plane of the orbit, that is, orthogonal to the initial velocity
vector vi. According to the definition given above, a simple plane change modifies
the direction, but not the intensity, of the velocity vector. Therefore, the three
vectors vi, vf, and Dv form an isosceles triangle.
With reference to the preceding figure, by applying the sine rule to the two
halves of this isosceles triangle, there results

1 h
D v ¼ v sin
2 2
and consequently the increment of velocity Dv can be expressed as follows

h
D v ¼ 2v sin
2
The preceding figure also shows that the angle formed by the two vectors vi and
Dv is 90° + h/2.
4.14 Change of Orbital Plane 703
A plane-change manoeuvre, if aimed at changing only the orbital inclination,

must be performed at one of the two nodes of the line of intersection between the
initial orbital plane and the final orbital plane. This is because the nodes are the only
two points which are common to both orbits. This manoeuvre, when performed at
one of the nodes, modifies only the orbital inclination, and does not affect the right
ascension of the ascending node and the argument of perigee.
Of course, a change of plane can be performed at any point of an orbit. However,
by so doing, there results a change not only in the orbital inclination, i, with respect
to the equatorial plane but also in the right ascension, X, of the ascending node.
Sometimes, this change is necessary to prepare a spacecraft for a manoeuvre of
rendezvous with another spacecraft. In these cases, it is desired to make a plane
change in order for the orbits of the two vehicles to lie in the same plane.
With reference to the following figure, let us consider two orbits (Orbit 1 and
Orbit 2) which have the point P3 in common and also intersect the equatorial plane
in the points, respectively, P1 and P2.
Let i1 and i2 be the inclinations of the two orbits with respect to the equatorial
plane. Let X1 and X2 be the right ascensions of the ascending nodes of the two
orbits. Let u1 = x1 + /1 be the argument of latitude relating to the position of the
first vehicle at a given epoch t1.
By applying the cosine rule to the spherical triangle P1P2P3, there results
cos h ¼ cos i1 cos i2 þ sin i1 sin i2 cosðX2 X1 Þ

cos i2 þ cos h cos i1
cos u1 ¼
sin h sin i1
As an example of application, let us consider an artificial satellite which, at a

given epoch t1, revolves about the Earth in a circular orbit of radius r1 = 11480 km,
inclination i1 = 60°.0, and right ascension of the ascending node X1 = 10°.0.
We want to determine the time and the difference of velocity which are required
to change the inclination of the orbit of this satellite to i2 = 45°.0 and the right
ascension of the ascending node to X2 = 55°.0.
The cosine of the angle, h, of the plane change required results from
cos h ¼ cos i1 cos i2 þ sin i1 sin i2 cosðX2 X1 Þ ¼ cos 60

:0 cos 45
:0
þ sin 60
:0 sin 45
:0 cosð55
0 10
:0Þ ¼ 0:78657
Hence, the angle of plane change is
h ¼ arccosð0:78657Þ ¼ 38
:134
In order to determine the time of the manoeuvre, we compute the cosine of the
argument of latitude u1 = x1 + /1 as follows
cos i2 þ cos h cos i1 cos 45

:0 þ cos 38
:134 cos 60
:0
cos u1 ¼ ¼
sin h sin i1 sin 38
:134 sin 60
:0
¼ 0:58684
Hence, the argument of latitude of the satellite at epoch t1 is
u1 ¼ arccosð0:58684Þ ¼ 125
:93
Since the initial orbit is circular, the argument of perigee may be set arbitrarily to
zero (x1 = 0°.0). Therefore, the variation of velocity must be performed when the
true anomaly of the satellite is
/1 ¼ 125
:93
In order to compute the difference of velocity required for the manoeuvre, it is

necessary to compute first the velocity, v, of the satellite in its orbit, as follows
12 1
lE 398600 2
v¼ ¼ ¼ 5:8925 km/s
r1 11480
Hence, the required difference of velocity results from

h 38 :134
Dv ¼ 2v sin ¼ 2 5:8925 sin ¼ 3:8498 km/s
2 2
As has been shown in this example, a simple plane-change manoeuvre is very

expensive in terms of energy consumption. As Braeunig [67] points out, if the angle
h of plane change were equal to 60°, the increment of velocity Dv, resulting from
the equation given above, would be equal to the orbital velocity v. In order to keep
this increment as small as possible, it is advisable to perform this manoeuvre where
the orbital velocity has its minimum value. In case of an elliptic orbit, this point is
of course the apogee.
Sometimes, it is necessary to change not only the inclination of the orbital plane
with respect to the equator but also the size of an initial orbit. Such is the case with a
satellite, which has been first launched from a site located at a nonzero geographic
latitude (that is, above or below the equator) and then inserted into a low circular
orbit of radius r1 with respect to the centre of mass the Earth. This satellite is to be
inserted into a geostationary orbit. The transfer manoeuvre necessary for this pur-
pose may be performed in two steps. In the first step, the first impulse of a
Hohmann transfer manoeuvre inserts the satellite into an elliptic orbit having its
perigee at the radius r1 of the initial circular orbit and its apogee at the radius
r2 = 42164 km. The second impulse of the same manoeuvre makes the orbit cir-
cular of radius r2. In the second step, a third impulse, relating to a simple
plane-change manoeuvre, reduces the inclination of the high circular orbit to zero,
that is, makes the orbit equatorial.
It is also possible to perform simultaneously the tangential burn at the apogee of
the elliptic orbit and the manoeuvre for plane change. By so doing, it is necessary to
change at the same time the magnitude and the direction of the velocity vector of a
spacecraft. For this purpose, let vi and vf (with |vi| 6¼ |vf|) be the velocity vectors of
the spacecraft at the apogee of the elliptic orbit, respectively, before a combined
plane-change-and-circularisation manoeuvre and after the combined manoeuvre.
Since the two vectors vi and vf differ not only in direction but also in magnitude,
then the variation of velocity at apogee must be such that
12
Dv ¼ jvi j2 þ jvf j2 2jvi jjvf j cos h
For example, let us consider an artificial satellite revolving about the Earth in an
orbit of radius r1 = 6678 km whose plane is inclined of i = 28°.5 with respect to
the equatorial plane. It is desired to move this satellite to a circular equatorial orbit
of radius r2 = 42164 km.
When the total manoeuvre is performed in two steps, a Hohmann transfer
manoeuvre moves the satellite to a circular orbit of radius r2 = 42164 km and
inclined of 28°.5 with respect to the equator. As has been shown in Sect. 4.12, the
total difference of velocity required by a Hohmann transfer manoeuvre is
ðDvÞH ¼ 3:8926 km=s

As has also been shown, the velocity in a circular orbit of radius r2 = 42164 km
is
12 1
lE 398600 2
v¼ ¼ ¼ 3:0747 km/s
r2 42164
When the satellite crosses the equatorial plane, a simple plane-change

manoeuvre is performed. This manoeuvre requires a further difference of velocity

h 28 :5
ðDvÞPC ¼ 2v sin ¼ 2 3:0747 sin ¼ 1:5137 km/s
2 2
Dv ¼ ðDvÞH þ ðDvÞPC ¼ 3:8926 þ 1:5137 ¼ 5:4063 km=s
By contrast, when a combined manoeuvre is performed at the apogee (that is,

when the satellite is at a distance r2 from the centre of the Earth) of the transfer
semi-ellipse whose major semi-axis is aH = ½(r1 + r2), then the initial velocity (in
km/s) of the satellite is
12 12
2 2 2 2
jvi j ¼ lE ¼ 398600 ¼ 1:6078
r2 r1 þ r2 42164 6678 þ 42164
and the final velocity of the satellite is
jvf j ¼ 3:0747 km=s
Therefore, the partial difference of velocity due to the combined manoeuvre is

12
ðDvÞC ¼ jvi j2 þ jvf j2 2jvi jjvf j cos h ¼ ð1:60782 þ 3:07472
1:6078 3:0747 cos 28

:5Þ2 ¼ 1:8303 km/s
1
This partial difference of velocity is to be added to the difference of velocity

necessary to insert the satellite from a circular orbit of radius r1 = 6678 km into a
transfer semi-ellipse having a perigee radius r1 = 6678 km and a major semi-axis
aH = (6678 + 42164)/2 km.
As has been shown in Sect. 4.12, this difference of velocity is
12 1
2 2 398600 2
ðDvÞ1 ¼ 398600 ¼ 2:4258 km/s
6678 6678 þ 42164 6678
Therefore, the total difference of velocity due to the combined manoeuvre is
Dv ¼ ðDvÞ1 þ ðDvÞC ¼ 2:4258 þ 1:8303 ¼ 4:2581 km=s
This value, compared with the value (Dv = 5.4063 km/s) relating to the two-step
manoeuvre, shows the advantage which can be obtained by performing a combined
manoeuvre. This advantage is due to the low value (1.6078 km/s) of the velocity of
the spacecraft at the apogee of the transfer semi-ellipse. By contrast, in case of a
two-step manoeuvre, a simple plane change is performed at one of the two nodes,
when the value (3.0747 km/s) of the spacecraft velocity in its circular orbit is higher
than the previous one.
Another possible option, suggested by Braeunig [67], consists in performing the
manoeuvre by means of three impulses. The first impulse, acting in the plane of the
low circular orbit, inserts the satellite into a semi-elliptic transfer orbit, whose
apogee is placed at a distance much greater than the radius of the geosynchronous
orbit. When the satellite reaches the apogee of this semi-ellipse, a combined
plane-change manoeuvre is performed. This manoeuvre inserts the satellite into
another semi-ellipse, which lies in the plane of the final circular orbit and whose
perigee radius coincides with the radius of the high circular orbit. Finally, when the
satellite reaches the perigee of the second semi-ellipse, a third impulse, acting in the
plane of the high circular orbit, inserts the satellite into this orbit. However, with
reference to the example considered above, since 42164 < 11.94 6678, then it is
impossible to spend less propulsive energy in a LEO-to-GEO transfer by using a
three-impulse combined manoeuvre instead of two-impulse combined manoeuvre.
In addition, a three-impulse manoeuvre takes more time than does a two-impulse
manoeuvre.
By the way, a spacecraft, which is placed at a distance r from the centre of mass
of a main attracting body and moves at a velocity v with respect to an inertial
system of reference whose origin is in the centre of mass of that body, has a
mechanical energy per unit mass
l 1 2
E ¼ þ v
r 2
where l is the gravitational parameter of the main attracting body. When the
velocity of that spacecraft increases instantaneously from v to v + Dv as a result of
an impulsive manoeuvre, then the mechanical energy per unit mass of the spacecraft
becomes
l 1
Eþ ¼ þ ðv þ DvÞ2
r 2
The increase of mechanical energy per unit mass
1
DE ¼ E þ E ¼ v D v þ ðDvÞ2
2
resulting from the instantaneous variation Dv is to be provided to the spacecraft by

its propulsion system.
4.15 Change of the Position of a Spacecraft in Its Orbit
The position (indicated by the true anomaly) of a spacecraft in its orbit can be
changed by means of a Hohmann transfer manoeuvre, which first takes the
spacecraft away from its original orbit, and then, after the lapse of a given time,
takes it back to that orbit in the desired position. The Hohmann transfer ellipse,
which is used for this purpose, is also called the phasing orbit (shown in the
following figure), whose period is chosen by the designer in such a way as to make
it possible to the spacecraft to return to its original orbit within a specified time.
With reference to the following figure, let /A be the true anomaly at a given time of
an artificial satellite which revolves about the Earth in a circular orbit of radius r1.
Let /B be the desired true anomaly of that satellite in the same orbit. It is desired to
move the satellite from A to B, that is, to change its true anomaly from /A to /B, so
as to reduce the difference D/ = /B − /A to zero after the lapse of a given time.
To this end, the satellite is moved from its circular orbit to a phasing elliptic orbit
(subscript PH) whose perigee radius is r1 and whose period TPH results from
imposing the following condition
4.15 Change of the Position of a Spacecraft in Its Orbit 709
nðkTPH Þ ¼ k ð2pÞ þ D/
where n = (lE/r31)½ is the mean motion of the satellite in its circular orbit, and k is
an integral number of revolutions which the satellite performs in the phasing ellipse
before returning to its circular orbit. The values of r1 and D/ (expressed in radians)
are supposed to be known. The value of k is chosen by the designer.
The preceding equation n(kTPH) = k(2p) + D/, solved for TPH, yields
2k p þ D/
TPH ¼
kn
The major semi-axis, aPH, of the phasing ellipse results from the well-known
equation
3 12
a
TPH ¼ 2p PH
lE
This equation, solved for aPH, yields

" 2 #13
TPH
aPH ¼ lE
2p
Since the velocity v = (lE/r1)½ of the satellite in its circular orbit is known, then
the difference of velocity, Dv1, which is required to insert the satellite in the phasing
ellipse of major semi-axis aPH, can be computed as has been shown in Sect. 4.12,
that is,
12 12
2 1 l
D v1 ¼ lE E
r1 aPH r1
Since this phasing ellipse is a Hohmann ellipse, its radius of apogee, r2, results
from the following equation
2aPH ¼ r1 þ r2
This equation, solved for r2, yields
r2 ¼ 2aPH r1
At the apogee of the phasing ellipse, another difference of velocity, Dv2, is

applied so as to insert again the satellite into its original orbit. The difference of
velocity Dv2 is equal in magnitude and opposite in sign to Dv1. Therefore, the total
difference of velocity required for a change of position is Dv = 2|Dv1|.
As an example of application, let us consider a geostationary satellite, whose

circular orbit is seen from the northern hemisphere of the Earth, as shown in the
preceding figure.
From this point of view, the satellite moves in the counterclockwise direction.
As is well known, the radius of the geostationary orbit is r1 = 42164 km. Let A be
the position of the satellite at a given time, and let /A be the true anomaly of the
satellite at that time. It is desired to shift the longitude of the satellite by an angle
D/ = 15° westward after the lapse of an interval of time corresponding to three
revolutions of the satellite in an elliptic phasing orbit, in order for the satellite to
occupy the position B at the final time.
The mean motion corresponding to the geostationary orbit is
12 1
lE 398600 2
n¼ ¼ ¼ 7:2922 105 rad/s
r13 421643
By imposing the condition
nðkTPH Þ ¼ k ð2pÞ þ D/
where, in the present example, k = 3 and D/ = 15p/180 rad, the orbital period of
the required phasing ellipse can be computed as follows
3:1416
2k p þ D/ 2 3 3:1416 þ 15 180
TPH ¼ ¼ ¼ 87361 s
kn 3 7:2922 105
The corresponding major semi-axis of this phasing ellipse results from

13 1
2
lE TPH 398600 873612 3
aPH ¼ ¼ ¼ 42554 km
4p2 4 3:14162
The difference of velocity, Dv1, necessary at the beginning of the manoeuvre,

that is, at the perigee of the phasing ellipse, is
12 12 12
2 1 l 2 1
Dv1 ¼ lE E ¼ 398600
r1 aPH r1 42164 42554
12
398600
¼ 0:014040 km/s ¼ 14:040 m/s
42164
The difference of velocity, Dv2, necessary at the end of the manoeuvre, that is, at
the apogee of the phasing ellipse, is equal in magnitude and opposite in sign to Dv1,
that is,
4.15 Change of the Position of a Spacecraft in Its Orbit 711
Dv2 ¼ Dv1 ¼ 14:040 m=s
The total difference of velocity required for the manoeuvre is then
Dv ¼ Dv1 þ jDv1 j ¼ 2 14:040 ¼ 28:080 m=s
4.16 Change of the Apsidal Line of an Orbit
The manoeuvre described in this paragraph has the purpose of transferring a

spacecraft from an elliptic orbit to another elliptic orbit having one of its two foci in
the same point as the first but its apsidal line rotated by an angle h with respect to
the apsidal line of the first orbit. An opportunity for performing such a transfer by
means of a single impulsive manoeuvre occurs in each point of intersection between
the two orbits. With reference to the following figure, let h be the angle formed by
the apsidal lines of two elliptic orbits having a common focus in the centre of mass,
F, of the Earth.
Let P and Q the points of intersection between the two orbits. The point P has a
common radius vector, r, but two different true anomalies, respectively, /1 and /2,
for the two different orbits. Each true anomaly is measured from the position of
perigee. The angle h is equal to the difference of the two true anomalies, as follows
h ¼ /1 /2
When the elements of the two orbits and the angle h are known, then the points
of intersection can be determined as will be shown below.
The radius vector r FP results from

p1 p2
r¼ ¼
1 þ e1 cos /1 1 þ e2 cos /2
where p1 = h21/lE, p2 = h22/lE, h1 and h2 are the moments of momentum per unit
mass of the spacecraft in the two elliptic orbits, and lE is the gravitational parameter
of the Earth.
The preceding equation may be written as follows
h22 e1 cos /1 h21 e2 cos /2 ¼ h21 h22
Since /2 = /1 − h, then
cos /2 ¼ cosð/1 hÞ ¼ cos /1 cos h þ sin /1 sin h
Substituting the preceding expression into h22 e1 cos /1 − h21 e2 cos /2 = h21 − h22
yields
A cos /1 þ B sin /1 ¼ C
where A = h22e1 − h21e2 cos h, B = − h21e2 sin h, and C = h21 − h22.

The equation A cos /1 + B sin /1 = C, whose coefficients A, B, and C are
known, can be solved as follows. By dividing the terms on both sides by A (which
coefficient is equal to zero only when the two ellipses are identical), there results
B C
cos /1 þ sin /1 ¼
A A
After setting
B sin a
¼
A cos a
the preceding equation can be written as follows
C
cos /1 cos a þ sin /1 sin a ¼ cos a
A
that is,
C
cos½ð/1 aÞ ¼ cos a
A
4.16 Change of the Apsidal Line of an Orbit 713
The ambiguity of sign is due to the two points of intersection between the two
ellipses. Hence, the general solution of the equation A cos /1 + B sin /1 = C is

C
/1 ¼ a arccos cos a
A
where

B
a ¼ arctan
A
Since /1 has been determined as has been shown above, then /2 is also
determined, because /2 = /1 − h.
As an example of application of the method described above, with reference to
the preceding figure, let rP1 = 10000 km and rA1 = 20000 km be the radius of
perigee and the radius of apogee of an elliptic orbit (Orbit 1). Let rP2 = 9000 km
and rA2 = 27000 km be the radius of perigee and the radius of apogee of another
elliptic orbit (Orbit 2), whose apsidal line is rotated counterclockwise by an angle
h = 30° with respect to the apsidal line of the first orbit (Orbit 1). We want to
compute the true anomaly, /1, and the difference of velocity, Dv, necessary for a
spacecraft to move from the first orbit to the second orbit at a point P.
The eccentricities, e1 and e2, of the two orbits are respectively
rA1 rP1 20000 10000 1

e1 ¼ ¼ ¼ ¼ 0:33333
rA1 þ rP1 20000 þ 10000 3
rA2 rP2 27000 9000 1
e2 ¼ ¼ ¼ ¼ 0:5
rA2 þ rP2 27000 þ 9000 2
The major semi-axes, a1 and a2, of the two orbits are respectively
rA1 þ rP1 20000 þ 10000

a1 ¼ ¼ ¼ 15000 km
2 2
rA2 þ rP2 27000 þ 9000
a2 ¼ ¼ ¼ 18000 km
2 2
The squares of the moments of momentum per unit mass are respectively

h21 ¼ a1 1 e21 lE ¼ 15000 1 0:333332 398600 ¼ 5:3147 109 km4 =s2

h22 ¼ a2 1 e22 lE ¼ 18000 1 0:52 398600 ¼ 5:3811 109 km4 =s2
The coefficients A, B, and C of the equation A cos /1 + B sin /1 = C are
A ¼ h22 e1 h21 e2 cos h ¼ 5:3811 109 0:33333 5:3147 109 0:5 cos 30
¼ 0:50762 109 km4 =s2

B ¼ h21 e2 sin h ¼ 5:3147 109 0:5 sin 30
¼ 1:3287 109 km4 =s2
C ¼ h21 h22 ¼ 5:3147 109 5:3811 109 ¼ 0:066433 109 km4 =s2
The angles a and /1 are

B 1:3287
a ¼ arctan ¼ arctan ¼ 69
:091
A 0:50762

C
0:066433
/1 ¼ a arccos cos a ¼ 69 :091 arccos cos 69 091

A 0:50762
¼ 69
:091 87
:323
The true anomaly of the point P in the first orbit is
/1 ¼ 69
:091 þ 87
:323 ¼ 156
:41
The true anomaly of the other point Q in the first orbit is
/1 ¼ 69
:091 87
:323 þ 360
¼ 341
:77
The true anomaly of the point P in the second orbit is
/2 ¼ /1 h ¼ 156
:41 30
:00 ¼ 126
:41
The radius vector r FP is

a1 1 e21 15000 ð1 0:333332 Þ
r¼ ¼ ¼ 19198 km
1 þ e1 cos /1 1 þ 0:33333 cos 156
:41
The tangent of the flight path angle, c1, in the point P of the first orbit is
e1 sin /1 0:33333 sin 156

:41
tan c1 ¼ ¼ ¼ 0:19207
1 þ e1 cos /1 1 þ 0:33333 cos 156
:41
Hence
c1 ¼ arctanð0:19207Þ ¼ 10
:872
The magnitude of the velocity vector of the spacecraft in P (in the first orbit) is
1 12
2 1 2 2 1
jv1 j ¼ lE ¼ 398600 ¼ 3:8668 km/s
r a1 19198 15000
The tangent of the flight path angle, c2, in the point P of the second orbit is
e2 sin /2 0:5 sin 126

:41
tan c2 ¼ ¼ ¼ 0:57222
1 þ e2 cos /2 1 þ 0:5 cos 126
:41
Hence
c2 ¼ arctanð0:57222Þ ¼ 29
:779
The magnitude of the velocity vector of the spacecraft in P (in the second orbit)
is
1 12
2 1 2 2 1
jv2 j ¼ lE ¼ 398600 ¼ 4:4024 km/s
r a2 19198 18000
The difference of velocity necessary for the manoeuvre is

h i12
j2Dvj ¼ jv1 j2 þ jv2 j2 2jv1 jjv2 j cosðc2 c1 Þ
After substituting the values computed above into the preceding equation, we
find
1
jD vj ¼ 3:86682 þ 4:40242 2 3:8668 4:0424 cosð29
:779 10
:872Þ 2
¼ 1:4573 km/s
By inspection of the preceding figure, the angle, d, formed by Dv with respect to

tho local horizontal results from
jv2 j sin c2 jv1 j sin c1 4:4024 sin 29

:779 3:8668 sin 10
:872
tan d ¼ ¼
jv2 j cos c2 jv1 j cos c1 4:4024 cos 29
:779 3:8668 cos 10
:872
¼ 61:602
By local horizontal, we mean the straight line produced from P perpendicularly

to FP, where P is the point of intersection indicated above, and F is the focus which
the two ellipses have in common. Hence
d ¼ arctanð61:602Þ ¼ 89
:070
By contrast, when the impulsive manoeuvre which changes the direction of the
apsidal line is performed in a given point P (of given true anomaly /1) of the first
orbit, then it is necessary to determine the elements and the angle of rotation, h, of
the second orbit. A possible way, indicated by Curtis [76], is described below. The
impulsive manoeuvre performed in P changes the radial component, vr, and the
transverse component, v/, of the velocity vector, v, of the spacecraft.
The moment of momentum per unit mass resulting from this change of velocity
h2 ¼ rðv/ þ Dv/ Þ ¼ h1 þ rDv/
where Dv/ is the transverse component of the vector Dv.

The radial component, vr, of the velocity vector v can be written as follows
lE
vr ¼ e sin /
h
This is because
dr h2 1 d/
vr ¼ r¼ h ¼ r2
dt lE 1 þ e cos / dt
Now, the expression vr = (lE/h) e sin / is applied to the second ellipse (sub-
script 2) at point P. At this point, the radial velocity and the true anomaly of the
spacecraft are respectively
vr2 ¼ vr1 þ Dvr

/2 ¼ /1 h
By so doing, there results

lE
vr2 ¼ vr1 þ Dvr ¼ e2 sin /2
h2
Substituting h2 = h1 + rDv/ into the preceding equation yields

lE lE
vr2 ¼ vr1 þ Dvr ¼ e1 sin /1 þ Dvr ¼ e2 sin /2
h1 h1 þ r D v /
The preceding equation, solved for sin /2, yields
l e1 sin /1 þ h1 Dvr
sin /2 ¼ h1 þ r D v/ E
e2 lE h1
On the other hand, the radius vector, r, relating to the common point P is
2
h21 1 h22 1 h1 þ r D v / 1
r¼ ¼ ¼
lE 1 þ e1 cos /1 lE 1 þ e2 cos /2 lE 1 þ e2 cos /2
The preceding equation, solved for cos /2, yields

2
h1 þ r D v/ e1 cos /1 þ 2h1 þ r D v/ r D v/
cos /2 ¼
e2 h21
Dividing the two preceding expressions of sin /2 and cos /2 yields
h1 lE e1 sin /1 þ h1 D vr
tan /2 ¼ h1 þ r D v / 2
lE h1 þ r D v/ e1 cos /1 þ 2h1 þ r D v/ r D v/
By replacing lE e1 sin /1 with h1 vr1, and h1 with r v/1, the preceding expression

rv2/1 v/1 þ D v/ ðvr1 þ D vr Þ
tan /2 ¼
lE v/1 þ D v/ 2 e1 cos /1 þ 2v/1 þ D v/ D v/
As shown in the preceding equation, the angle, h, formed by the two apsidal
lines of the two orbits is determined by the two components, Dvr and Dv/, of the
impulse vector Dv, which is applied to the spacecraft at the point in which the true
anomaly is /1. When these two components are specified, then the preceding
equation makes it possible to compute tan /2 and hence /2 = arctan(tan /2). Then,
the preceding expressions of sin /2 or cos /2 make it possible to compute the
eccentricity, e2, of the second orbit. Finally, h2 results from h2 = h1 + rDv/.
The equations indicated above hold for the general case of an impulse vector
(Dv) having arbitrary radial (Dvr) and transverse (Dv/) components, which vector is
applied to a spacecraft at any point (identified by the true anomaly /1) of its orbit.
In the particular case of an impulse applied to a spacecraft at the perigee (where

/1 = 0 and vr1 = 0) of its initial orbit and only in the radial direction (Dv/ = 0),
then the preceding equation becomes
rv/1 D vr
tan /2 ¼ tanðhÞ ¼
l E e1
Because of the minus sign in front of h, when the radial impulse is applied to a
spacecraft in the outward direction (Dvr > 0), then the apsidal line of the second
orbit rotates clockwise with respect to the apsidal line of the first orbit; in the same
manner, when the radial impulse is applied to the spacecraft in the inward direction
(Dvr < 0), then this rotation is counterclockwise.
As an example of application, let rP1 = 10000 km and rA1 = 20000 km be,
respectively, the radius of perigee and the radius of apogee of an elliptic orbit (Orbit
1). Let Dv = 1.5 km/s be the magnitude of an impulse applied to a spacecraft,
which revolves about the Earth along the orbit indicated above, at the perigee (rP1)
of this orbit. Let c1 = 60° be the angle which the impulse vector Dv forms with
respect to the local horizontal. We want to determine the new orbit (Orbit 2)
followed by the spacecraft after this manoeuvre.
The eccentricity, e1, of the first orbit is
rA1 rP1 20000 10000 1

e1 ¼ ¼ ¼ ¼ 0:33333
rA1 þ rP1 20000 þ 10000 3
The major semi-axis, a1, of the first orbit is
rA1 þ rP1 20000 þ 10000

a1 ¼ ¼ ¼ 15000 km
2 2
The velocity of the spacecraft at the perigee of the first orbit has only the
transverse component, v/1, which results from the vis-viva integral
12 12
2 1 2 1
v/1 ¼ lE ¼ 398600 ¼ 7:2901 km/s
rP1 a1 10000 15000
The radial component and the transverse component of the impulse vector are
respectively
Dvr ¼ Dv sin c1 ¼ 1:5 sin 60

¼ 1:2990 km=s
Dv/ ¼ Dv cos c1 ¼ 1:5 cos 60
¼ 0:75 km=s
The true anomaly of the spacecraft in the second orbit, after the application of
the impulse vector, is

rv2/1 v/1 þ D v/ ðvr1 þ D vr Þ
tan /2 ¼ 2
lE v/1 þ D v/ e1 cos /1 þ 2v/1 þ D v/ D v/
where r = rP1 = 10000 km, v/1 = 7.2901 km/s, lE = 398600 km3/s2,

Dv/ = 0.75 km/s, vr1 = 0, Dvr = 1.2990 km/s, e1 = 0.33333, and cos /1 = 1.
By substituting these values in the preceding equation, we have
tan /2 ¼ 0:42140
Hence
/2 ¼ arctanð0:42140Þ ¼ 22
:85
Since in the present case /2 = –h and Dvr > 0, then the apsidal line rotates
clockwise, after the application of the impulse vector, through an angle of 22°.85.
The eccentricity, e2, of the second orbit results from the equation given above
2
h1 þ r D v/ e1 cos /1 þ 2h1 þ r D v/ r D v/
cos /2 ¼
e2 h21
where h1 = r v/1 and cos /1 = 1. By solving this equation for e2, we have
2
v/1 þ D v/ e1 þ 2v/1 þ D v/ D v/
e2 ¼
v2/1 cos /2
ð7:2901 þ 0:75Þ2 0:33333 þ ð2 7:2901 þ 0:75Þ 0:75
¼ ¼ 0:67474
7:29012 cos 22
:85
The moment of momentum per unit mass of the spacecraft in the second orbit
results from
h2 ¼ h1 þ r D v/ ¼ rv/1 þ r D v/ ¼ 10000 ð7:2901 þ 0:75Þ ¼ 80401 km2 =s
The major semi-axis, a2, of the second orbit results from
h22 804012
a2 ¼ ¼ ¼ 29772 km
lE 1 e22 398600 ð1 0:674742 Þ
The radius of perigee and the radius of apogee of the second orbit result from
rP2 ¼ a2 ð1 e2 Þ ¼ 29772 ð1 0:67474Þ ¼ 9683:6 km

rA2 ¼ a2 ð1 þ e2 Þ ¼ 29772 ð1 þ 0:67474Þ ¼ 49860 km
4.17 Drag Make-up Manoeuvres for Satellites

in Low-Altitude Orbits
As has been shown in Sects. 3.19 and 3.20, the principal effect produced by the
aerodynamic drag on a satellite revolving around the Earth at low altitudes is a
decrease in both the major semi-axis and the eccentricity of the satellite orbit. Drag
make-up manoeuvres are meant to maintain the major semi-axis of a satellite orbit
within an acceptable range of tolerance, depending on the mission which has been
planned for that satellite.
The first of the Lagrange planetary equations in Gaussian form, which concerns
specifically the rate of change with time of the major semi-axis, may be written as
follows
2
da 2a v
a0 ¼ aDt
dt lE
where a is the major semi-axis, v is the magnitude of the orbital velocity of a

satellite, aDt is the tangential (i.e. in the direction of the velocity vector) component
of the perturbing acceleration vector aD due to drag acting on the satellite, and lE is
the gravitational parameter of the Earth.
The preceding equation may also be written as follows
2
2a v
Da ¼ Dvt
lE
where Da is the increment of major semi-axis produced by an impulse Dvt applied

to the satellite in the direction of the velocity vector at a given point of its orbit.
The preceding equation shows that the best point for executing a manoeuvre of
correction of the major semi-axis of an orbit is the perigee of that orbit, because the
orbital velocity reaches there its maximum value.
The amount of decay, Da, in major semi-axis due to aerodynamic drag within a
given interval of time can be computed as has been shown in Sect. 3.20. For
example, a first-approximation value for the amount of decay within a given
interval of time, Dt, can be computed by using a Taylor-series expansion truncated
after the first-order term, as follows

a20
Da B q0 v30 Dt
lE
where B = CDA/m is the ballistic coefficient of the satellite.

By the way, it is to be noted that the components (aDr, aD/, and aDh) of the
perturbing acceleration vector aD due to drag have been determined and expressed
in Sect. 3.20 in the radial (subscript r), transverse (subscript /) and bi-normal
(subscript h) co-ordinates. By contrast, when it is desired to express the same vector
4.17 Drag Make-up Manoeuvres for Satellites in Low-Altitude Orbits 721
in the tangential (subscript t), normal (subscript n), and bi-normal (subscript h)
co-ordinates, then the following transformation can be used to compute aDt and aDn
from aDr and aD/
" # " #
aDt vp2 r 2 ðe sin /Þ=p 1=r aDr
¼ 2 2 2
aDn h r e sin / þ p2 1=r ðe sin /Þ=p aD/
As a result of this transformation, the equation of Sect. 3.2
2a2 h p i
a0 ¼ ðe sin /ÞaDr þ aD/
h r

2
2a v
a0 ¼ aDt
lE
When Da has been determined, the velocity impulse, Dvt, which is to be applied
to the satellite at perigee in the direction of the velocity vector in order to com-
pensate for this decay, results from the equation Da = (2a2v/lE)Dvt solved for Dvt,
as follows
l
Dvt ¼ E
Da
2a2 v
When this velocity impulse has been determined, the principle of conservation of
momentum can be used to estimate the mass of propellant required for this
manoeuvre. For this purpose, it is necessary to know the mass of the satellite and
the specific impulse of its thruster.
As has been shown in Sect. 4.5, the specific impulse, Isp, is proportional to the
velocity, ve, at which the exhaust gases leave the final section of the nozzle
ve
Isp ¼
g0
where g0 is the acceleration of gravity at the surface of the Earth. In Sect. 3.25, it
has also been shown that the change in mass, Dm, of the satellite due to the thrust
can be expressed by means of the Tsiolkovsky rocket equation

D v
D m m1 m2 ¼ m1 1 exp
g0 Isp
where m1 results from summing the structural mass (ms), which in turn comprises
the structure, the engines, the propellant tanks, the control systems, et c., the
payload mass (mpay), and the propellant mass (mprop), as follows
m1 ¼ ms þ mpay þ mprop
and m2 is the mass at burnout time. An example will be given in Sect. 4.19.
In particular, when the difference (Dv) of the satellite velocity due to the thrust
can be neglected in comparison with the product g0Isp, then exp[−Dv/(g0Isp)] can be
expanded in a Maclaurin series truncated after the first-degree term. For this pur-
pose, remembering that
1 2 1 3
expðxÞ ¼ 1 x þ x x þ 1 x
2 6
the following expression can be used for a rough estimate of the mass Dm of
propellant required for the manoeuvre
Dv
Dm m1
g0 Isp
4.18 Manoeuvres for Geostationary Satellites
The basic manoeuvres to be performed in order to insert an artificial satellite from a

low-altitude (rP 6678 km) orbit around the Earth into a geostationary orbit have
been shown in the preceding paragraphs. In the present paragraph, the matter is
considered in detail.
After inserting a satellite into a low-altitude orbit, it is necessary to move it to a
geosynchronous transfer orbit, that is, to a high-eccentricity orbit, whose apogee is
located in the equatorial plane at a radius rA = rGEO = 42164 km and whose
inclination with respect to the equatorial plane is equal to or greater than the
geocentric latitude of the launch site, depending on the safety constraints discussed
in Sect. 4.4. In particular, according to the table given in Sect. 4.4, the launch site of
Cape Canaveral, Florida, is located at a latitude of 28°.5N, and consequently the
minimum inclination possible for a geosynchronous transfer orbit is 28°.5 in case of
a launch due East. The launch site of Kourou, French Guyana, is located at a
latitude of 5°.2N; however, due to different mission constraints, the inclination
obtained for a geosynchronous transfer orbit is usually 7° [77].
A geosynchronous transfer orbit has usually its apogee and its perigee in the
equatorial plane. When a satellite has been inserted into this orbit, it is necessary to:
(a) reduce the high eccentricity of this orbit to zero, and (b) reduce the inclination of
the resulting circular orbit to zero. As has been shown in Sect. 4.14, the final result
may be obtained by performing either two separate manoeuvres or one combined
manoeuvre. It has also been shown in the same paragraph that a single combined
plane-change-and-circularisation manoeuvre is advantageous over a two-step
manoeuvre, because, in case of a launch from Cape Canaveral, it requires a total
4.18 Manoeuvres for Geostationary Satellites 723
expenditure, expressed in terms of Dv, of 4.2581 km/s (comprising 2.4258 km/s for
the transfer orbit and 1.8303 km/s for the combined manoeuvre at apogee) against
5.4063 km/s (comprising 2.4258 km/s for the transfer orbit, 1.4668 km/s for the
circularisation at apogee, and 1.5137 km/s for the simple plane change at a node).
The corresponding expenditure in terms of mass of propellant may be computed by
means of the equation

D v
Dm m1 m2 ¼ m1 1 exp
g0 Isp
given in Sect. 4.17, in which the initial mass, m1, and the specific impulse, Isp, are
also to be specified. In practice, as Sidi [77] points out, the combined manoeuvre
requires more than one motor firing. This is because the orbit actually obtained after
applying the impulse Dv to the satellite may differ from the desired geostationary
orbit as a result of attitude errors. In order to correct these errors with the minimum
expenditure of propellant, the total impulse Dv is divided into two or more
impulses, as follows
Dv ¼ a Dv þ ð1 aÞDv
The preceding equation refers to the case of two impulses, where a is a number,
to be determined, falling between zero and unity. The criterion for determining the
value of a is the minimum expenditure of propellant necessary to correct an angular
error, a, committed during the manoeuvre. This error has a magnitude of a few
degrees.
With reference to the following figure, let AB vi and AC vf be the vectors
indicating, respectively, the initial velocity of the satellite at the apogee of the
geosynchronous transfer orbit and the final velocity of the satellite necessary for the
desired geostationary orbit.
In the presence of an angular error a, if the vector impulse BD = Dv were not

divided into two impulses BE and EC, then another impulse DC of magnitude
a
2D v sin
2
(the sides BD and BC of the triangle BCD have the same length Dv) would be
necessary in order for the satellite to achieve the desired velocity vector vf. As
shown by the preceding figure, this would imply a high expense of propellant.
For the purpose of reducing the expense of propellant to a minimum in case of
an angular error a committed in the firing stage, the impulse vector Dv is divided
into two parts, namely aDv and (1 − a)Dv, where the unknown coefficient a must
be determined. With reference to the preceding figure, since a is a small angle, then
tan a a, and consequently there results
EF h ¼ ðaDvÞa
h i12 1
BE ¼ ðaDvÞ2 þ h2 ¼ aDv 1 þ a2 2
n o12 h i12
EC ¼ ½ð1 aÞDv2 þ h2 ¼ Dv ð1 aÞ2 þ a2 a2
1
¼ Dv a2 1 þ a2 þ 1 2a 2
n 1 1 o
BE þ EC ¼ D v a 1 þ a2 2 þ a2 1 þ a2 þ 1 2a 2
As Sidi [77] points out, there are two velocity terms, denoted by ev1 and ev2,
which determine the further Dv due to the attitude error a committed in the firing
stage. The first term is
n 1 1 o
v1 ¼ BE þ EC ðBF þ FCÞ ¼ Dv a 1 þ a2 2 þ a2 1 þ a2 þ 1 2a 2 D v
and the second term is

1
v2 ¼ ðECÞa ¼ Dva a2 1 þ a2 þ 1 2a 2
The sum of these two velocity terms is

n 1 1 o
v ¼ v1 þ v2 ¼ Dv a 1 þ a2 2 þ ð1 þ aÞ a2 1 þ a2 þ 1 2a 2 1
The unknown quantity, whose value can reduce the sum ev to a minimum, is
a. By differentiating the preceding expression of ev with respect to a, and equating
the result to zero, there results
dðvÞ 1 2að1 þ a2 Þ 2
¼ 1 þ a2 2 þ ð 1 þ aÞ 1 ¼ 0
da 2½a2 ð1 þ a2 Þ þ 1 2a2
After setting for convenience A = 1 + a2 and B = (1 + a)2, the preceding

equation can be written as follows
dðvÞ 1 1 aA 1
¼ A2 þ B2 1 ¼ 0
da ða A þ 1 2aÞ2
2
Therefore, the first derivative of ev with respect to a vanishes when
1 1 1
A2 a2 A þ 1 2a 2 ¼ B2 ð1 aAÞ
By squaring the terms on both sides of the preceding equation, there results

A a2 A þ 1 2a ¼ Bð1 aAÞ2
ð1 BÞA2 a2 þ 2AðB 1Þa þ A B ¼ 0
By solving the preceding equation for a, discarding the plus sign in front of the
square root, and taking account of the settings A = 1 + a2 and B = (1 + a)2, there
results the value of the coefficient a which leads to the minimum expenditure of
propellant in case of an angular error a committed in the firing stage, as follows
" 12 #
1 a
a¼ 1
1 þ a2 aþ2
where the value of the angular error a is expressed in radians. The following table
shows the values of the coefficient a computed as a function of the angular error a
(expressed in degrees).
a (°) a a (°) a
0.0000 1.0000 5.0000 0.78952
1.0000 0.90671 6.0000 0.76851
2.0000 0.86797 7.0000 0.74888
3.0000 0.83798 8.0000 0.73031
4.0000 0.81239 9.0000 0.71257
The ratio of the velocity term ev = ev1 + ev2 to the impulse Dv, expressed in per
cent, results from
v n 1 1 o
ð%Þ ¼ a 1 þ a2 2 þ ð1 þ aÞ a2 1 þ a2 þ 1 2a 2 1 100
Dv
The value of this ratio makes it possible to compute the expenditure of propellant
due to an angular error a.
As has been shown in Sect. 3.13, the principal causes which perturb the orbit of
a geostationary satellite are:
• the luni-solar gravitational attraction, which causes a secular variation of the

inclination vector (see below);
• the solar radiation pressure, which moves the eccentricity vector (see below)
along a circle as the Sun moves, as seen from the Earth, in a near-circular orbit;
and
• the non-spherical shape of the Earth, which moves the satellite in longitude.
The principal orbital elements affected by such perturbations are the inclination
(i) of the orbital plane with respect to the equatorial plane, the eccentricity (e) of the
orbit, and the longitude (k) of the satellite, which changes with respect to the
desired longitude (k0). In order to counter-act the perturbations indicated above, the
following manoeuvres are necessary:
• longitude (east–west) station-keeping;
• inclination (north–south) station-keeping; and
• eccentricity corrections.
These manoeuvres are considered below.
The problem of east–west station-keeping for geostationary satellites has been
studied by several authors, for example, Kamel et al. [78], Romero et al. [79], Wie
and Roithmayr [80], and Lee et al. [81]. Since the satellite orbit is circular and
equatorial, in order to avoid numerical indetermination, the following set of orbital
elements is frequently used: the drift rate (D) in mean longitude, the two compo-
nents (iX = i sin X and iY = −i cos X) of the inclination vector i, the two compo-
nents (eX = e cos(X + x) and eY = e sin(X + x)) of the eccentricity vector e, and
the sidereal angle (s = hG + k = hG0 + xE(t − t0) + k) where hG is the Greenwich
sidereal angle, and k is the sub-satellite longitude.
The eccentricity vector, e, defined in Sect. 1.1, goes from the centre of mass, O,
of the Earth to the perigee of the orbit.
The inclination vector, i, shown above, is the projection OP of the unit vector
rv
uh ¼
j r vj
onto the equatorial plane XY, and is perpendicular to the line of nodes. The X-axis
and the Y-axis shown in the preceding figure lie in the equatorial plane of the Earth,
where the origin O is the centre of mass of the Earth, and the X-axis points towards
the first point of Aries. The bi-dimensional inclination vector, i, is represented by
the segment OP, whose length is sin i i, in case of an orbit of small inclination
with respect to the equator.
The segment OP results from projecting the unit vector uh = (r v)/|r v| onto
the equatorial plane perpendicularly to the line of nodes. Hence, the inclination
vector, i, has magnitude sin i i, and points in the direction X − 90° [82]. The
drift rate is defined as follows
3 Da
D¼
2 aGEO
where aGEO = 42164 km is the major semi-axis of the geostationary orbit, and
Da = a − aGEO is the deviation of the major semi-axis from this value [82].
When the six orbital elements indicated above are known, the east–west (k),
north–south (u), and radial (q) relative co-ordinates of the satellite result from
u ¼ iX cos s iY sin s

3
q ¼ aGEO aGEO D þ eX cos s þ eY sin s
2
k ¼ k0 þ Dðs s0 Þ þ 2eX sin s 2eY cos s
where the subscript 0 refers to t = t0. As to the east–west station-keeping, we take

account here only of the perturbation due to elliptical shape of the Earth equator,
which is the principal cause of the drift in longitude. In these conditions, by con-
sidering only the first (C22 and S22) of the tesseral harmonic coefficients and the
case of a geostationary orbit (geocentric latitude equal to zero, and r = a), the radial
(ar), transverse (a/), and bi-normal (ah) components (see Sect. 3.2) of the per-
turbing acceleration can be expressed as follows
9lE rE2
ar ¼ ½C22 cosð2kÞ þ S22 sinð2kÞ
a4
6lE rE2
a/ ¼ ½C22 cosð2kÞ S22 sinð2kÞ
a4
ah ¼ 0
where rE is the mean equatorial radius of the Earth and a = 42164 km is the major
semi-axis, or the radius, of the circular geostationary orbit. Since

0 da 2
a ¼ a/
dt n
where n = (lE/a3)½ is the mean motion of the orbit, then the perturbing acceleration
in longitude (k″ d2k/dt2) can be computed as follows
da 3 1 5 3 1
dn d 12 3 d 1

k00 ¼ l2E a2 ¼ l2E a2 2a2 lE 2 a/
3
¼ lE a 2 ¼
dt dt da dt 2

3
¼ a/
a
The preceding expression a/ = –(6lEr2E/a4)[C22 cos(2k) − S22 sin(2k)], substi-

tuted into this equation, leads to
18lE rE2 18lE rE2

k00 ¼ 5
½C 22 cosð2kÞ S22 sinð2kÞ ¼ J 22 sin½2ðk k22 Þ
a a5
where J22 = (C222 + S222)½, and k22 = −0.26056 rad = −14°.929 −15°.0.
There are four equilibrium longitudes, denoted by k*, separated by approxi-
mately 90° along the equator. These equilibrium longitudes, for which k″ = 0,
correspond to the following condition
S22 0:90387 106

tanð2k Þ ¼ ¼ ¼ 0:57405
C 22 1:5745 106
Hence
k ¼ 75
; 165
; 255
; and 345
It can be shown that the longitudes k* = 75° and k* = 255° are points of stable
equilibrium, and the other two (k* = 165° and k* = 345°) are points of unstable
equilibrium. Since the two longitudes of stable equilibrium (denoted by ks) are
separated by ±p/2 from k22, then
18lE rE2 18lE rE2 n h p io

k00 ¼ 5
J 22 sin½2ðk k22 Þ ¼ 5
J 22 sin 2 k ks
a a 2
18lE rE2
¼ J 22 sin½2ðk ks Þ¼ 0:00168 sin½2ðk ks Þð
=day2 Þ
a5
where ks are the two longitudes of stable equilibrium indicated above. Agrawal [83]
has shown that, when higher tesseral harmonic coefficients than C22 and S22 are
taken into account, the expression of the perturbing acceleration in longitude k″
differs from that given above, and that the points of stable equilibrium are located at
the following longitudes ks = 73°.9 and ks = 256°.3.
In practice, a geostationary satellite is placed in its orbit within a narrow (±0°.1

or even ±0°.05) range, also called dead band, which is centred around the desired
longitude k0. In these conditions, the perturbing acceleration k″ can be considered
independent of longitude and constant. By so doing, the preceding differential
equation k″ = −0.00168 sin[2(k − ks)] can be readily integrated as follows
1
kðtÞ ¼ k0 þ k00 ðt t0 Þ þ k00 ðt t0 Þ2
2
where t0 is the initial time, and k0 = k(t0), k′0 = k′(t0). The function k(t) indicated
above can be represented in the phase plane (k, k′) as a parabola within the limits
k0 − Dk and k0 + Dk, where Dk is the designed semi-width of the narrow interval
of longitude, as shown in the following figure. In this figure, the perturbing
acceleration in longitude (k″) is assumed to be negative (k″ < 0), and therefore the
parabola, which describes the drift cycle of the satellite around k0, is concave
towards the left-hand side of the phase plane (k, k′). In the case of positive per-
turbing acceleration in longitude (k″ > 0), the reverse is true.
Let us suppose that the satellite is in the point of the phase plane identified by the
abscissa k0 − Dk and the ordinate k′0 at the initial time t0. In order to keep the
satellite within the desired longitude interval, it is necessary to impart an impulse
Dv to the satellite when it reaches the western limit of the desired interval at a drift
velocity equal to −k′0. As a result of this impulse, the satellite comes back to the
previous point (k0 − Dk, k′0) of the phase plane.
It is necessary to determine the drift velocity for a given perturbing acceleration

|k″| and a desired semi-width Dk. This done, it is necessary to compute the time
interval s between two consecutive east–west station-keeping manoeuvres and the
impulse Dv to be applied to the satellite in each manoeuvre.
In the case (shown in the preceding figure) of a negative perturbing acceleration

(k″ = –|k″|) of constant value, by integrating successively the differential equation
k″ = –|k″|, there results
k0 ðtÞ ¼ jk00 jt þ k00

1
kðtÞ ¼ jk00 jt2 þ k00 t þ k0
2
Let t1 be the time such that k′(t1) = 0. As a result of this definition, there results
k0 ðt1 Þ ¼ jk00 jt1 þ k00 ¼ 0
and hence
k00
t1 ¼
jk00 j
The preceding figure shows that t1 is the time necessary for the longitude of the
satellite to change from k0 −Dk to k0 + Dk, and therefore t1 is half the total time, tf,
necessary for the complete cycle of drift. Hence there results
k00
tf ¼ 2t1 ¼ 2
jk00 j
During the time t1, the longitude of the satellite changes from k0 −Dk to
k0 + Dk, and therefore
0 2 0
1 1 k k0
kðt1 Þ ¼ jk00 jt12 þ k00 t1 þ k0 ¼ jk00 j 000 þ k00 Dk
2 2 jk j jk00 j
2
1 k00
¼ Dk ¼ Dk
2 jk00 j
The preceding equation, solved for k′0, yields the required drift velocity
1
k00 ¼ 2ðjk00 jDkÞ2
Hence, the drift velocity for a given perturbing acceleration |k″| is determined.
Let s be the time interval required by the drift velocity to the value change from
the value +k′0 to –k′0. The preceding equation k′(t) = –|k″| t + k′0 yields
k0 ðsÞ ¼ jk00 js þ k00 ¼ k00

The preceding equation, solved for s, yields
k00
s¼2
jk00 j
As has been shown above, this is also the total time, tf, necessary for the
complete cycle of drift. Hence, by introducing the expression k′0 = 2(|k″|Dk)½ into
the equation written above, there results
1
Dk 2
s¼4
jk00 j
Finally, the impulse, Dv, which is to be imparted to the satellite in order to

change its drift velocity from +k′0 to –k′0, can be computed as will be shown below.
Firstly, the energy integral
1 2 lE l
v ¼ E
2 r 2a
is differentiated keeping the radius vector, r, constant, because only the velocity of a
spacecraft is affected, by definition, in an impulsive manoeuvre. This differentiation
yields
lE lE D a
vDv ¼ 2
2Da ¼
ð2aÞ 2a2
The preceding expression, solved for Da, yields

2vDva2 2a
Da ¼ ¼ Dv
lE v
Secondly, the equation (third law of Kepler) which defines the mean motion
l 12 1
¼ l2E a2
3
n¼ E
a3
is differentiated. This yields
3 1 5
Dn ¼ l2 a2 Da
2
Substituting Da = 2vDva2/lE into the preceding equation yields

3vDv
Dn ¼ 1
ðlE aÞ2
The preceding equation, solved for Dv, yields

1
1 ðlE aÞ2 1 na2
Dv ¼ Dn ¼ Dn
3 v 3 v
which expresses the impulse to be imparted to the satellite in order to change its
drift velocity from +k′0 to –k′0.
In case of a geostationary satellite, there results
a ¼ 42164 km
l 12 39860012
n ¼ E3 ¼ ¼ 7:2921 105 rad/s
a 421643
l 12 39860012
v¼ E ¼ ¼ 3:0747 km/s
a 42164
By introducing these values into the equation
1 na2
Dv ¼ Dn
3 v
derived above, there results

7:2921 105 421642
Dv ¼ Dn ¼ 14054Dn km=s
3 3:0747
The change in mean motion, Dn, is evaluated by Sidi [77] by means of the
change in drift velocity, Dk′, which in turn is equal to 2k′0 for each manoeuvre, as
shown in the preceding figure.
By measuring the mean motion, n, in degrees/day (instead of radians/second)
and the impulse, Dv, in m/s (instead of km/s), the preceding equation becomes
14054 1000
Dv ¼ ¼ 2:8391Dk0 m=s
180
24 60 60
3:1416
Since, for each manoeuvre, Dk′ = 2k′0, then, for each manoeuvre, Dv (in m/s) is
h 1
i h 1
i
D v ¼ 2:8391 2k00 ¼ 5:6782k00 ¼ 5:6782 2ðjk0 jDkÞ2 ¼ 11:357 ðjk00 jDkÞ2
Since the time interval between two consecutive manoeuvres has been deter-
mined above (s = 4(Dk/|k″|)½), then the impulse (in m/s) to be imparted to the
satellite in a year can be computed as follows
2 3
6 7
1 365 1
6 365 7
D v ¼ 11:357 ðjk00 jDkÞ2 ¼ 11:357 ðjk00 jDkÞ2 6 1 7
s 4 Dk 2 5
4
jk00 j
¼ 1036:3 jk00 j ¼ 1036:3j0:0017 sin½2ðk ks Þj ¼ 1:7618 sin½2ðk ks Þ
The preceding equation shows that the impulse necessary to keep a satellite
within the desired interval of longitude in a given time does not depend on the
amplitude (Dk) of the interval. It depends on the difference in longitude (k − ks)
between the satellite and a point of stable equilibrium.
As has been shown above, the corrective manoeuvre for the east–west station
keeping of a geostationary satellite is to be performed in the direction opposite to
the drift in longitude. For example, let us consider a geostationary satellite designed
for a longitude k0 = 120°E, which is greater than the longitude, ks 75°E, of a
point of stable equilibrium. That satellite tends to drift westward to the stable point
of longitude ks, due to the tesseral harmonics of the gravitational field of the Earth.
The westward drift increases the major semi-axis and hence the period of the orbit.
In this case, the required east–west station-keeping manoeuvre should be performed
to decrease the major semi-axis of the satellite

2a
Da ¼ Dv
v
For this purpose, thrusters placed on the east side of the satellite should be used
for negative transverse impulses Dv [81].
In case of other manoeuvres than that for the east–west station-keeping, the
impulse vector, Dv, to be applied to the satellite has three components: a radial
component, Dvr, a transverse component, Dv/, and a bi-normal component, Dvh.
Soop [82] expresses these components in terms of the six orbital elements D, eX, eY,
iX, iY, and s as follows
Dvr ¼ vðDeX sin sb DeY cos sb Þ

Dv/ ¼ vðDD þ 2DeX cos sb þ 2DeY sin sb Þ
Dvh ¼ vðDiX sin sb DiY cos sb Þ
where the subscript b indicates that the sidereal angle, s, refers to the burn time of
the thrusters. The effects produced by these components, for the purpose of com-
pensating for the deviations caused by the remaining perturbations (solar radiation
pressure and luni-solar gravitational force), are considered below.
The solar radiation pressure is another cause of drift in the instantaneous plane of
motion. This pressure, which is particularly important in a satellite having large
solar panels, induces an eccentricity into an orbit and leaves the other orbital
elements unchanged [84]. In other words, the solar radiation pressure causes
variations in the eccentricity vector e. This is because, as the solar radiation pressure
accelerates the satellite on one side of the orbit, the orbital radius decreases on the
other side. There, the solar radiation pressure slows the satellite, so that the orbit
becomes still more elongated [85]. The magnitude of this effect depends on the
ballistic coefficient (see Sect. 3.19) of the satellite.
Chu et al. [85] give the following expression for the average rate of change in the
eccentricity vector, integrated over a day:

de 3pS A sin sS
¼
dt 2vm cos sS
where pS is the solar radiation pressure (see Sect. 3.16), A is the area of the cross
section of the satellite, v and m are, respectively, the velocity and the mass of the
satellite, and sS is the sidereal angle of the Sun in the equatorial plane, as shown in
the following figure.
According Chu et al., for ballistic coefficients greater than 50 kg/m2, eccentricity
increases less than 0.4 10−6 per day and can be controlled by choosing the
manoeuvre time [85], as will be shown below.
For a nearly-circular orbit, as is the case with geostationary orbits, the change,
De, in the eccentricity vector is the result of an impulse, Dv, applied to the satellite
in the instantaneous plane of motion. Let Dvr and Dv/ be, respectively, the radial
component and the transverse component of the impulse Dv. The resulting changes
in the eccentricity vector are

2 1
DeX ¼ cos s Dv/ þ sin s Dvr
v v

2 1
D eY ¼ sin s D v/ þ cos s D vr
v v
(References [77] and [84]), where v = 3.0747 km/s is the velocity of the satellite in
the geostationary orbit, and s is the right ascension of the satellite at the burn time
(that is, s sb). The method indicated by Pocha and Edwards for controlling the
oscillations in longitude induced by solar radiation pressure consists in controlling
the eccentricity of the orbit and the position of perigee with respect to the Earth–
Sun line, as shown in the preceding figure. Initial conditions are set up in order for
the orbit to have some eccentricity, with the position of perigee at some angle
behind the Sun line. The eccentricity will then decrease, until it reaches a minimum
value half-way through the cycle, at which point the perigee is located along the
Sun line. At the end of the cycle, the eccentricity will reach again its initial value,
and the perigee will be located in front of the Sun line by the same amount as it was
behind that line at the beginning of the cycle. At this point, a manoeuvre is per-
formed in order to regain the initial conditions. This manoeuvre moves the perigee
position, while the eccentricity remains unchanged [84].
According to Agrawal [83, page 85], the eccentricity is generally corrected
indirectly, by performing a longitude correction at apogee by means of a transverse
(Dv/ only) impulse in the direction of motion in order to increase the major
semi-axis, followed by another transverse (Dv/ only) impulse at the perigee of the
resulting orbit, oppositely to the direction of motion, in order to decrease the major
semi-axis.
The resultant solar radiation pressure acts on the centre of pressure of a satellite.
When this point does not coincide with the centre of mass, then that satellite is
subject to a torque. In case of an axially symmetric satellite having a single solar
array, the primary component of this torque is in the orbital plane and is perpen-
dicular to the Sun line. In order to compute the increase in the moment of
momentum h = r mv of the satellite due to solar radiation pressure, it is nec-
essary to integrate Euler’s equations (see Sect. 7.6) in co-ordinates x, y, and z ro-
tating with the satellite at angular velocity x
M ¼ h0 ¼ h þ x h
(Refence [85]), where M is the torque due to solar radiation pressure, h′ dh/dt,
and h is the total moment-of-momentum vector of the satellite, resulting from the
sum of the part, = x, due to the rigid body and the part, hw, due to the control
wheels
h ¼ = x þ hw
For a perfect system of attitude control with no angular accelerations, the dif-
ferential equation for the part of moment of momentum due to the control wheels is
[85]:
hw ¼ M x ð=
= x þ hw Þ
The moment-of-momentum vector rotates in the satellite body co-ordinates,

because it is tied to the Sun line.
The inclination of the orbit of a geostationary satellite is subject to change as a
result of the gravitational forces exerted by the Sun and by the Moon. A change in
the orbit inclination causes a motion in latitude (north–south) of the sub-satellite
point. The manoeuvre intended to counteract this effect is called north–south
station-keeping. As is the case with the chosen longitude, k0, of the sub-satellite
point, there is also a range, sometimes called dead band, within which the (geo-
centric) latitude, u, is allowed to vary around the required value u0 = 0, corre-
sponding to an equatorial orbit.
The semi-length, Du0, of the latitude range is usually taken equal to ± 0.1 or
even ±0.05°. A north–south station-keeping manoeuvre, which is meant to correct
the inclination of a satellite orbit with respect to the equator, is usually described in
terms of the motion of the pole of the orbital plane of the satellite.
By pole of an orbital plane, we mean the point where the line of action of the
unit vector
rv
uh ¼
j r vj
intersects the spherical surface of unit radius whose centre is the centre of mass of
the Earth. In other terms, with reference to the following figure, the spherical
co-ordinates (radius, r, right ascension, a, and declination, d) of the orbital pole are
r¼1 a ¼ X 90
d ¼ 90
i
The orbital pole defined above has, in the geocentric equatorial system XYZ, the
following co-ordinates
X ¼ sin i sin X i sin X

Y ¼ sin i cos X i cos X
Z ¼ cos i
Let P be the projection of the orbital pole onto the equatorial plane. If all
perturbations were absent, the abscissa X and the ordinate Y of P would not change
with time. Instead, as a result of the luni-solar perturbation, X and Y are functions of
time. Therefore, latitude corrections are necessary in order to counteract the effects
of the luni-solar perturbation and keep the satellite within the designed dead band.
The variations of X and Y include periodic and secular terms. The rates of change,
di/dt and dX/dt, in orbital inclination and in the right ascension of the ascending
node due to the secular terms have been computed by Agrawal and can be found in
Ref. [83].
The manoeuvre performed for latitude station-keeping must compensate for the
variations induced by the secular terms, so as to reduce the expenditure of pro-
pellant to a minimum, and must also be compatible with the manoeuvre performed
for longitude station-keeping.
The method indicated by Pocha and Edwards [84] consists in setting up the
initial conditions for a north–south station-keeping manoeuvre in order for the orbit
inclination at the initial time and at the final time of a cycle to be the same.
Let s be the start time of a cycle, XS and YS be the abscissa and the ordinate of P
at the start time of a cycle, XE and YE be the abscissa and the ordinate of P at the end
time of a cycle, sC be the duration of a cycle, X′ and Y′ be the rates of change of
X and Y due to the secular terms, and DX(t) and DY(t) be the cyclic offsets of X and
Y from secular motion at time t due to the secular terms of the luni-solar
perturbation.
Since we want the orbital inclination to be the same at the start time and at the
end time of a cycle, then the inclination vector, i, must have the same magnitude at
these two times, that is, the following equality must hold
XS2 þ YS2 ¼ XE2 þ YE2
The ordinate of P is set to zero at the start time of a cycle. In other words, we set
YS ¼ DYðsÞ
Therefore, the ordinate of P at the end time of a cycle results
YE ¼ Y 0 sC þ DYðs þ sC Þ
The abscissa of P at the end time of a cycle results
XE ¼ XS DXðsÞ þ X 0 sC þ DXðs þ sC Þ
Now, the following expressions YS = DY(s), YE = Y′sC + DY(s + sC), and

XE = XS − DX(s) + X′sC + DX(s + sC) are substituted into X2S + Y2S = X2E + Y2E,
and this equation is solved for XS.
This yields
1 A
XS ¼ B þ
2 B
where
A ¼ ½DYðsÞ2 ½DYðs þ sC Þ2 Y 0 sC ½Y 0 sC þ 2DYðs þ sC Þ

B ¼ 2½DXðs þ sC Þ DXðsÞ þ X 0 sC
In other words, by setting XS = −½ B + A/B and YS = DY(s), the orbit incli-

nation at the initial time and at the final time of a cycle are the same. The abscissa
and the ordinate of the point P change over a cycle, of duration sC, from (XS, YS) to
(XE, YE).
A north–south station-keeping manoeuvre can be designed so that, after a further
period of time of duration sC, the abscissa should change from XE to XS(s + sC),
and the ordinate should change from YE to YS(s + sC). Pocha and Edwards [84]
have determined the impulse, Dvh, necessary to obtain this result as follows
n o12
Dvh ¼ v ½XS ðs þ sC Þ XE 2 þ ½YS ðs þ sC Þ YE 2
where v = 3.0747 km/s is the velocity of a satellite in a geostationary orbit.

They have also determined the right ascension, aN or aS, at which this impulse is
to be applied to the satellite as follows

X S ð s þ sC Þ X E
aN ¼ arctan
Y S ð s þ s C Þ YE

X S ð s þ sC Þ X E
aS ¼ arctan þp
Y S ð s þ s C Þ YE
where aN refers to northerly thrusts, and aS refers to southerly thrusts. By using this
method, Pocha and Edwards [84] found the following requirement for the L-SAT
(whose latitude dead band was Du0 = ±0.05°): 357 m/s for 7 years. The duration
of each cycle was sC = 16 days, and the maximum inclination over the 7-year
lifetime was 0.036°. The necessary impulse was found to vary from a maximum of
2.77 m/s to a minimum of 0.91 m/s. On the other hand, the requirement relating to
the east–west station-keeping for the same satellite (whose longitude dead band was
Dk0 = ± 0.087°) was found to be 36 m/s for 7 years, and hence the total
requirement was 357 + 36 = 393 m/s over the same life time.
4.19 De-orbiting Manoeuvres
Some space missions require the return of an orbiting vehicle to the surface of the
Earth. Such is the case with all manned missions, whose crew must be made to
return safely to some intended place of the Earth. Such is also the case with an
orbiting weapon designed to come out of its orbit in order to strike some place of
the Earth, by means of one or more impulses provided by a propulsion system and
meant to reduce the orbital velocity of that weapon [86].
A de-orbiting manoeuvre is a sequence of events which occur under human
control. By contrast, the re-entry of a space vehicle concerns the natural decay of a
spacecraft, previously launched from some place on the surface of the Earth, into
the lower layers of the atmosphere, due to the action of aerodynamic forces, such as
lift and drag. In these layers, the increasing density of the atmosphere affects
severely the trajectory of the spacecraft, and the aerodynamic heating burns up in
most cases the spacecraft, so as to prevent it from striking the Earth. In other cases,
some parts of large space structures can survive the passage through the atmosphere
and hit the Earth, which event poses a problem of safety, because the impact may
occur in populated regions of the Earth. In addition, there is a problem arising from
the increasing number of debris located especially at altitudes less than 2000 km
and at geostationary altitudes (35786 km). Such debris can cause serious damages
in case of collision with other spacecraft. For this reason, NASA has issued a safety
standard [87], according to which all current and future space missions must
consider the end of life, or the disposal, of spacecraft during their design, con-
struction, and operational employment. Another standard on the same matter has
been issued by the ESA [88]. As Novin [89] points out, there are several options for
a satellite which comes to the end of its life. A satellite orbiting at low altitudes may
be raised, at the end of its operational life, to a storage orbit of altitude greater than
2000 km, in order to remove it from its previous orbit, and insert it into another
orbit far away from other operational satellites. The same satellite may also be
lowered to a decaying orbit, in order to be destroyed by aerodynamic heating,
provided that there is a chance whose value is less than 1 out of 10000 of causing
injury to people on the ground. Still a further option is the actively de-orbiting the
satellite, so as to make its re-entry faster and manage its impact point throughout the
process. The third of these options is considered below.
According to Shidemantle [90], two requirements have to be taken into account in
de-orbiting a satellite. The first requirement is to keep the amount of the necessary
impulse (Dv) as small as possible, in order to decrease the expenditure of propellant
and hence the total cost of the mission. The second requirement is to keep the flight
path angle as high in absolute value as possible while the satellite plunges deeper
into the atmosphere, in order to decrease the chances of pieces breaking off the
satellite and skipping out of the atmosphere. In addition, a high absolute value of the
flight path angle results in a smaller impact footprint, because the pieces are less
scattered over long distances than would be the case otherwise. It is to be noted that
the flight path angle has negative values during the descent trajectory, because the
velocity vector has a component opposite to the position vector [91].
The Compton Gamma Ray Observatory (CGRO) was the first satellite de-orbited
intentionally by NASA. After its launch on the 5th of April 1991 aboard the Space
Shuttle Atlantis, the CGRO was safely de-orbited and entered the Earth atmosphere
on the 4th of June 2000. For the controlled crash, four de-orbit burns were performed
over a duration of 30 min by means of the hydrazine propulsion system. The first
burn lowered the orbit perigee altitude from 510 to 350 km. The second burn
occurred about 25 h later and lowered the perigee altitude to 250 km. After about
four days, the third burn lowered the perigee altitude to 150 km. The fourth burn
occurred on the next orbit and caused the CGRO to go deeper into the atmosphere,
where it began to tumble and heat up before being disintegrated. The impact foot-
print occurred about 4000 km southeast of Hawaii and was estimated to be 26 km
wide and 1552 km long. No pieces of debris struck the land [90]. The sequence of
burns is shown in the following figure, which is due to the courtesy of NASA [92].
4.19 De-orbiting Manoeuvres 741
Several burns were performed instead of one to decrease the perigee radius of the
CGRO satellite for the following reasons. Firstly, spreading the whole de-orbiting
manoeuvre over several days makes it possible to determinate accurately the orbits.
In case of errors committed in a burn, it is possible to correct them by a new
planning of the subsequent burns. This new planning occurs usually by operating
together with the Cheyenne Mountain Operations Centre. This way of proceeding
leads to a high accuracy in the determination of the impact area. Secondly, several
small burns produce less stresses and vibrations in the structure of the satellite than
those generated by a single burn of high intensity.
In case of the satellite to be de-orbited being part of a constellation of satellites, a
preliminary manoeuvre is necessary to remove the satellite of interest from its
operational orbit, in order to avoid any risk of collision with other satellites
belonging to the same constellation. Such was the case with the Spot 1 satellite,
whose original orbit was lowered in the preliminary phase 15 km below the original
orbit [93]. In the second phase, several burns were performed at the apogee of each
of the planned intermediate orbits, for the purpose of obtaining a sequence of
elliptic orbits of increasing eccentricity, whose apogee radius was constant and
whose perigee radius decreased progressively in the transition from one of the
intermediate orbits to the next. For this purpose, each burn was performed in the
direction opposite to the direction of the orbital velocity of the satellite at apogee.
The burn performed at the apogee of the last elliptic orbit is meant to reach the
radius of the final perigee. The satellite re-entry occurs in the same orbit as that of
the last burn. Of course, the target chosen should be an uninhabited region of the
Earth. The satellite is expected to break up at an altitude comprised between 84 and
70 km. In other words, within this interval, the satellite begins to disintegrate. Such
was the case with the CGRO satellite, as shown in the following figure, due to the
The de-orbiting manoeuvre described above, comprising several elliptic transfer

orbits meant to reduce progressively the perigee radius, is shown in the following
example. Let us consider the DMSP 5D-3/F17 (NSSDC/COSPAR ID: 2006-050A)
satellite, which was launched by a Delta IV vehicle from the Vandenberg AFB on
the 4th of November 2006 at 15:53:00 UTC. In the Celestial Observer [94], we find
the following information:
Satellite Elements of Object 2006-050A
Dimensions 3 m 2 m, cylindrical
Brightness 6.0 mag (at 1000 km and 50% illuminated)
4.9 mag (at perigee and full illumination)
NORAD catalog Nr 29522
International designator 2006-050A
Perigee height 837.6 kma (relative Earth’s equatorial radius)
Apogee height 852.7 kma (relative Earth’s equatorial radius)
Orbit period 101.8 mina
Kepler Elements
Major semi-axis 7223.274 kma

Eccentricity 0.00105
Inclination 98.8401°
Argument of perigee 286.9738°
Right ascension of ascending node 106.1985°
Mean anomaly at epoch 73.0261°
Susceptibility to drag Ba 3.7328e−05 1/Earth radii
Revolutions per day 14.141714a
Number of orbits since launch 0
Epoch (UT): 15 April 2010 16:06:26 (JD = 2455302.17114119)

State Vector
Position vector x = −2014.401 km

y = +6942.287 km
z = +14.339 km
Velocity vector vx = +1.0981 km/s
vy = +0.3112 km/s
vz = +7.3387 km/s
a
Original data calculated from SGP4 model elements
Two-Line Elements (TLE)
Since the orbital eccentricity is nearly zero, the orbit may be considered circular
with a radius, r0, equal to the major semi-axis indicated above, that is,
r0 = 7223.3 km. The velocity of the satellite in the initial circular orbit results from
v0 = (lE/r0)½.
In the present example, since DMSP F17 is part of a constellation comprising
two Sun-synchronous satellites, then we plan a preliminary phase during which the
radius of perigee is reduced by an amount of about 15 km, as suggested by Alby
[93]. Hence, for this phase, the perigee radius is chosen as follows
r1 ¼ r0 15:0 ¼ 7223:3 15:0 ¼ 7208:3 km
This phase is followed by operations of orbit determination. The radii of perigee

in the second phase depend on a choice to be made by of the designer. For a reason
which will be given below, we set a final radius of perigee r9 = 6417.9 km, that is,
a final altitude of perigee h9 = 39.81 km. In this phase, we choose eight decrements
of 98.8 km each, so that 15.0 + 8 98.8 + 6417.9 = 7223.3 km, and obtain the
following table, whose columns have been filled as will be shown below.
№ h (km) r (km) a (km) v (km/s) Dv (km/s) ½T (min)

0 845.21 7223.3 7.4285
1 830.21 7208.3 7215.8 7.4247 0.0038482 50.834
2 731.41 7109.5 7166.4 7.3990 0.025688 50.313
3 632.61 7010.7 7117.0 7.3728 0.026137 49.794
4 533.81 6911.9 7067.6 7.3462 0.026596 49.276
5 435.01 6813.1 7018.2 7.3192 0.027069 48.761
6 336.21 6714.3 6968.8 7.2916 0.027555 48.247
7 237.41 6615.5 6919.4 7.2636 0.028056 47.735
8 138.61 6516.7 6870.0 7.2350 0.028570 47.224
9 39.81 6417.9 6820.6 7.2059 0.029100 46.716
Totals 0.22262 438.90
The first or leftmost column contains an integral number which indicates the
orbit of reference, where the zero designates the initial circular orbit. The second
column contains the altitude (in km) of perigee for each orbit. The third column
contains the radius (in km) of perigee for each orbit. The fourth column contains the
major semi-axis (in km) for each transfer orbit. For the first transfer ellipse, we take
the value, r1 < r0, chosen for the perigee radius. Hence, the major semi-axis, a1, of
the first transfer ellipse results from
r0 þ r1
a1 ¼
2
Likewise, for the second transfer ellipse, we take the value, r2 < r1, chosen for
the perigee radius. Hence, the major semi-axis, a2, of the second transfer ellipse is
r0 þ r2
a2 ¼
2
and so on. The fifth column contains, in the row zero, the velocity v0 = (lE/r0)½ of
the satellite in the original circular orbit. The same column contains, in the rows
going from 1 to 9, the velocity of the satellite at the apogee of each transfer ellipse.
For example, the velocity of the satellite at the apogee of the first transfer ellipse
results from the vis-viva integral
12
2 1
v1 ¼ lE
r 0 a1
Likewise, the velocity at the apogee of the second transfer ellipse results from
1
2 1 2
v2 ¼ lE
r 0 a2
and so on. The sixth column contains the values of the impulses (Dv) to be applied
to the satellite at the common apogee of the transfer ellipses in order to decrease the
radius of perigee. For example, the impulse to be applied to the satellite at the
apogee of the first transfer ellipse, in order for the radius of perigee to decrease from
r0 = 7223.3 km to r1 = 7208.3 km, is
Dv1 ¼ v0 v1
Likewise, the impulse to be applied to the satellite at the apogee of the second
transfer ellipse, in order for the radius of perigee to decrease from r1 to r2, is
Dv2 ¼ v1 v2
and so on. The seventh column contains the semi-period of each transfer ellipse.
The period, Tk, of each transfer ellipse results from Kepler’s third law
12
a3k
Tk ¼ 2p
lE
where the subscript k = 1, 2, …, 9 indicates the transfer ellipse of interest. The

satellite takes a time equal to ½Tk = p(a3k /lE)½ to move from the common apogee to
the perigee of each transfer ellipse.
The choice of a radius of 6417.9 km for the perigee of the last transfer ellipse is
due to the requirement of a sufficiently high absolute value of the flight path angle
of the satellite at the entry point, ce, in order to make the impact area as small as
possible, and also to prevent the satellite from skipping out. By entry point we mean
the point of the orbit whose altitude is 120 km over the ground. Janovsky et al. [95]
recommend an altitude of perigee equal to or smaller than 60 km for the last ellipse,
and an absolute value of the flight path angle such that 1°.5 |ce| 2°.5,
depending on the altitude of the initial orbit, before the last braking burn. Ahmed
et al. [96] recommend an absolute value |ce| greater than 1.2 degrees. We choose
ce = −1.85 degrees. On the other hand, in the example considered above, the initial
circular orbit is a Sun-synchronous orbit, and the successive elliptic transfer orbits
are used to lower progressively the altitude of perigee. In such conditions, it seems
reasonable to perform the burns in the northern semi-orbit, when the satellite is
approximately over the north pole of the Earth. In addition, in accordance with
Novin [89], we choose a value of approximately 325° for the true anomaly, /e, of
the entry point, in order to place the satellite over the South Pacific en route with the
final perigee altitude and the impact zone. At any point of an elliptic orbit, the flight
path angle c of a satellite results from
e sin /
tan c ¼
1 þ e cos /
After solving this equation for e with ce = −1°.85 and /e = 325°, the eccen-
tricity, e9, of the final transfer ellipse results
tan ce
e9 ¼ ¼ 0:059036
sin /e cos /e tan ce
By introducing this value into the following equation
1 e9
r9 ¼ r0
1 þ e9
where r9 is the radius of perigee of the last transfer ellipse and r0 = 7223.3 km is
the common radius of apogee, there results
r9 ¼ 6417:9 km
which corresponds to an altitude of perigee h9 = 39.81 km. This explains the values
indicated in the row No. 9 of the preceding table. These values satisfy both of the
conditions 1°.5 |ce| 2°.5 and h9 60 km, which are given in Ref. [95].
The major semi-axis, a9, of the last ellipse is
r0 þ r9
a9 ¼ ¼ 6820:6 km
2
The radius vector, re, and the velocity, ve, of the satellite at the entry point are

a9 1 e29
re ¼ ¼ 6483:3 km
1 þ e9 cos /e
1
2 1 2
ve ¼ lE ¼ 8:0325 km/s
r e a9
In the last segment of the satellite trajectory, the acceleration due to aerodynamic
drag is the principal perturbation. As has been shown in Sects. 3.19 and 3.20, the
principal effects produced by drag are a decrease in major semi-axis and in
eccentricity. For a preliminary evaluation of these effects, first-approximation val-
ues of major semi-axis, a, and eccentricity, e, can be computed by using
Taylor-series expansions truncated after the first-order terms, as follows

a20
a a0 3
Bq0 v0 D t
l
E
ð1 e0 Þa0
e e0 3
Bq0 v0 D t
lE
where the subscript 0 indicates the values at the time of reference t = t0.
In the example considered above, in the absence of other data, the value of the
ballistic coefficient (B = CDA/m) of the DMSP 5D-3/F17 satellite can be computed
from the B* coefficient indicated in the two-line elements, as follows
B ¼ 12:741621B ¼ 12:741621 0:37328 104 ¼ 0:00047562 m2 =kg
However, this value is subject to change, because the mass of the satellite
decreases due to the consumption of propellant necessary for the burns, and also
because the solar panels of the satellite may be put in the folded position after the
last burn, as a result of a command coming from the control centre. In addition, the
tanks may be emptied of the residual propellant after the last burn, before the
electrical storage devices discharge, in order to avoid risks of explosion [97]. This is
because some propellants, as is the case with hydrazine, are highly reactive.
The table given above makes it possible to compute the mass of propellant
necessary for a de-orbiting manoeuvre, provided that the total mass (structure plus
propellant) of the satellite and the specific impulse of the thrusters are known, by
using the Tsiolkovsky rocket equation

D v
Dm m1 m2 ¼ m1 1 exp
g0 Isp
as has been shown in Sect. 4.17. For the example given above, the DMSP F17
satellite has a launch mass m1 = 1154 kg, including a dry mass of 825 kg [98].
Supposing the mass of the satellite before de-orbiting to be 1154 kg, and supposing
the thrusters to be the hydrazine thrusters CHT 400 N, whose specific impulse in
vacuo is Isp = 220 s [99], the amount of propellant necessary for the de-orbiting
manoeuvre may be computed as follows:
Dm1 ¼ 1154 f1 exp½3:8482=ð9:8 220Þg ¼ 2:0579 kg

Dm2 ¼ ð1154 2:0579Þ f1 exp½25:688=ð9:8 220Þg ¼ 13:644 kg
Dm3 ¼ 1138:3 f1 exp½26:137=ð9:8 220Þg ¼ 13:716 kg
Dm4 ¼ 1124:6 f1 exp½26:596=ð9:8 220Þg ¼ 13:787 kg
Dm5 ¼ 1110:8 f1 exp½27:069=ð9:8 220Þg ¼ 13:859 kg
Dm6 ¼ 1096:9 f1 exp½27:555=ð9:8 220Þg ¼ 13:930 kg
Dm7 ¼ 1083:0 f1 exp½28:056=ð9:8 220Þg ¼ 14:002 kg
Dm8 ¼ 1069:0 f1 exp½28:570=ð9:8 220Þg ¼ 14:072 kg
Dm9 ¼ 1054:9 f1 exp½29:100=ð9:8 220Þg ¼ 14:143 kg
Hence, the total amount of propellant necessary for the de-orbiting manoeuvre of
the DMSP 5D-3/F17 satellite is
DmT ¼ Dm1 þ Dm2 þ þ Dm9 113:21 kg
Determining the time and the area of impact of a satellite on the surface of the
Earth is a difficult task. As has been shown in Sect. 2.7 et seq., the values of many
terms which appear in the equations of motion are not known with certainty,
because the environment in which the satellite moves is not known exactly. For this
reason, the uncertainty which affects the data available to the designer is treated by
means of probabilistic methods. In other words, we try to determine the likelihood
of a particular impact area or time, because some physical quantities (e.g. heat
transfer coefficients, emissivity, temperature, density, etc.) which define the envi-
ronment cannot be known exactly. Among these methods, one of the most widely
used is the Markov Chain Monte Carlo method. The fundamental concepts on
which this method is based are briefly described below.
A Markov Chain is a series of trials arranged in a time sequence which satisfies
the following two conditions:
(a) the outcome of each trial belongs to a finite set of outcomes, which is called the
state space; and
(b) the outcome of each trial depends only on the outcome of the trial which
precedes immediately the trial considered.
A classical example of a Markov Chain is the symmetric random walk, also
known as drunkard’s walk. Let us imagine a drunk person, who can move only
leftward or rightward by one pace. In this case, the probability of a move on the left
side is equal to that of a move on the right side, and the future position of that
person depends only on the present position. Another example of a Markov Chain
is a board game like Monopoly, where the position occupied by a player after
rolling the die depends only on the position occupied by that player before the roll,
let the preceding positions be what they may. Still another example is the price of
the stocks of a company.
The Monte Carlo method, which is so called after the Monte Carlo casino, is a
computational method which uses random numbers. This is because there is an
analogy between what happens in a casino, where wins and losses occur by chance,
and the generation of random numbers. This method consists of calculating millions
of Markov Chains and analysing the results to get the best estimate. The results of
some mathematical or physical processes are too difficult or even impossible to be
predicted in a deterministic manner. In such cases, a computer simulation of the
process of interest provides the best estimate of the result unknown, because the
Markov chain converges to the desired probability distribution after a number of
steps. In the specific case of a satellite entering the atmosphere of the Earth, the
Markov Chain Monte Carlo method constructs a probability model based on the
observed data. This probability model is then used to estimate the value of any
statistical quantity of interest, for example, the impact time or the location of the
impact area. In mathematical terms, let h [h1, h2, …, hp]T be a vector containing
p quantities which are necessary for modelling the trajectory of a satellite. Let p(h)
be a prior probability distribution concerning the quantities h1, h2, …, hp. After the
collection of new data, the prior can be updated by means of the likelihood func-
tion, L(D|h), where D [d1, d2, …, dn]T is a vector containing the n new data
collected.
Following Horsley, the updating is performed according to Bayes’ theorem:
pðhÞLðDj hÞ
pðh jDÞ ¼ R R R
pðhÞLðDj hÞdh1 dh2 dhp
where p(h|D) is called the posterior distribution and represents the probability of the
p quantities h1, h2, …, hp after the collection of the new data. The expression in the
denominator of the preceding fraction implies the integration over all possible
satellite trajectories. In general, this multiple integral cannot be solved analytically.

In other words, the preceding equation cannot be solved exactly in many cases, and
in particular in those concerning the prediction of the trajectory of a satellite in the
atmosphere of the Earth.
An approximation to the true value of the multiple integral can be computed by
means of the Markov Chain Monte Carlo method. This method is used to generate a
large number of satellite trajectories, each of which is sampled in direct proportion
to its probability. In other words, such trajectories as are very likely to occur are
sampled many times by the procedure, and such trajectories as are not very likely to
occur are sampled a few times or are not sampled at all. By so doing, an ensemble
of trajectories is generated. Any statistical calculation, such as computing a mean,
variance or credible interval can be performed on this ensemble. Greater accuracy,
if needed, can be obtained by sampling a greater number of trajectories. This
method requires a large number of iterations. Since the numerical integration of the
equations of motion for a satellite travelling through the Earth atmosphere requires
a great amount of time, then this method implies the use of a large computer in
order to generate results in a timely manner.
This method begins by defining a Markov Chain, M, which is then used to
generate samples from a probability distribution p on the space of all possible
re-entry trajectories, X. This Markov Chain is constructed in such a way as to have
its states x belonging to the set X and a stationary distribution p(x).
Let x be the current state of the Markov Chain M. A new state, x′, is proposed
for the next step of the chain. The new state is generated following a proposal
distribution, q(x, x′), which depends on the current state of the chain.
The proposed new state (or move), x′, is accepted with a probability

pðx0 Þqðx0 ; xÞ
Aðx; x0 Þ ¼ min 1;
pðxÞqðx; x0 Þ
As a result of the ergodic theorem, it can be shown (see Refs. [100] and [101])
that a Markov Chain M converges to its stationary distribution as long as M
satisfies a few requirements, such as those of being irreducible and aperiodic.
By using a uniform proposal distribution, such as that, which is shown in the
algorithm given below, it is possible to simplify the acceptance probability as
follows

0 pðx0 Þ
Aðx; x Þ ¼ min 1;
pðxÞ
By so doing, the acceptance probability depends only on the ratio of the like-
lihoods and does not require the knowledge of the normalising constant Z.
MCMC algorithm for quantifying the uncertainty of re-entry. N(m, r) refers to a

random draw from a normal distribution having a mean of m and a variance of r.
(a) Initialise proposal (prop) trajectory by using current (curr) trajectory
(b) sample U from [1,2,3,4] uniformly at random
if U = 1 then (Area-to-mass ratio move)
A/mprop = N[(A/mcurr), (A/mcurr)/10.0]
else if U = 2 then (shape parameter 1 move)
s1prop = N(s1prop, s1prop/10.0)
end if
(c) compute trajectory and calculate likelihood
(d) accept/reject proposed re-entry trajectory by using A(x, x′)
(e) record resulting re-entry time and location.
By so doing, it is possible to generate a large number of re-entry trajectories in
direct proportion to their probability, but without the necessity of integrating over
all possible trajectories. It is necessary to compute the likelihood of a re-entry
trajectory. The logarithm of the likelihood function, L(D|h), results from computing
the chi-squared cost between the value of the major semi-axis, am, relating to the
model and the value of the major semi-axis, ad, relating to the data at the time of
validity of a set of two-line elements, and then summing over all the data points
used in the computation, as follows
X
N 2
1 aid aim
log½LðDj hÞ ¼
i¼1
2 r
where r is the variance of the date relating to the major semi-axis. Since the data
coming from the two-line elements have no estimate of the variance relating to the
major semi-axis, then the value of the variance must come from other sources.
Horsley [102] has found that a value of 100 metres yields good results.
A different acceleration model is used for each trajectory considered in the
Monte Carlo method. These models are based on all possible accelerations acting
upon a satellite. The sources of uncertainty in these models are the ballistic coef-
ficient, B, of the satellite considered, the atmospheric density, q, the solar proxy,
F10.7, and the geomagnetic activity index, ap, as has been shown in Sect. 3.22.
Horsley [102] used four parameters to model the uncertainty due to the ballistic
coefficient, and one parameter to model the uncertainty due to the atmospheric
density. He used a different stochastic model for each iteration to take account of
the uncertainty due to the solar activity. He did not consider the uncertainty due to
the geomagnetic activity. The algorithm used by Horsley for quantifying the
uncertainty of re-entry is given above. He estimated the probability of a particular
re-entry time by means of the following expression
1 Xnmc
Pbt ¼ Nðt\tn \t þ dtÞ
nmc nbi n¼nbi
where P b t is the probability of the satellite re-entering at time t, nmc is the number of
iterations performed in the Monte Carlo method, and nbi is an initial number of
discarded moves performed in the calculation.
By so doing, the probability is computed by counting the number, N, of re-entry
trajectories which occur between t and t + dt, and dividing this number by the total
number of iterations used in the calculation.
The complexity and the length of the calculations which have been shown above
make it necessary to use fast and accurate tools for predicting the footprint area and
the time of impact. It is also necessary to assess the risks due to surviving fragments
which may hit the ground and cause damage to people in populated regions. The
tools available for this purpose have been described at length by Lips and Fritsche
[103], and then by Wu et al. [104]. A brief account of the results found by these
authors is given below.
The methods, on which the tools available at present for this purpose are based,
have been classified by Lips and Fritsche into two categories. They are the
object-oriented method and the spacecraft-oriented method. The first of these
methods consists in simplifying the complex shapes of real satellites by means of
simple shapes, such as spheres, cylinders, boxes, etc. The principal tools based on
this method are
• Debris Assessment Software (DAS), developed by NASA;
• Object Reentry Survival Analysis Tool (ORSAT), also developed by NASA;
and
• Debris Reentry and Ablation Prediction System (DRAPS), developed by Wu
et al.
DAS and ORSAT are described in Ref. [105]. They are the computer codes used
by NASA for predicting the survivability of satellites and components coming from
the upper stages of launch vehicles which enter the atmosphere from either orbital
decay or controlled manoeuvres. This prediction is necessary to determine the risks
to people on the ground. After DAS has determined that the re-entry of a given
spacecraft does not comply with the NASA Technical Standard STD-8719.14 [87],
ORSAT is used for a deeper analysis of survivability. ORSAT uses trajectory,
atmospheric, aerodynamic, and thermal-ablation models for both a satellite and its
launch vehicle in order to evaluate the impact risk. The trajectory of a moving body
is considered in three degrees of freedom. The atmospheric models used in ORSAT
are either the 1976 U.S. Standard Atmosphere, or the MSIS-E-90 atmosphere, or
the GRAM-99 atmosphere. After defining the shape for the object of interest by
means of simpler shapes, ORSAT considers the possibility for the object entering
the atmosphere in either spinning mode or tumbling mode. The drag coefficient
attributed to this object is considered in a range going from hypersonic speed to
subsonic speed in order to evaluate the kinetic energy of the object at the moment of
its impact with the ground. The heating rates of spherical objects are adjusted for
other bodies at various regimes of rarefied flow. Both lumped mass and
one-degree-of-freedom models are used to compute the surface temperature. The
object is considered to demise when the heat absorbed by it reaches the heat of
ablation of the material of which it is made.
The preceding figure, due to the courtesy of NASA [105], illustrates the pre-
dicted break-up of the Upper Atmosphere Research Satellite (UARS) when entering
the atmosphere. ORSAT made it possible to predict demise altitude versus down-
range for most components of the UARS.
DRAPS is another computing tool based on the object-oriented method.
DRAPS, according to its authors [104], uses a three-degree-of-freedom ballistic
model for trajectory computation and either a lumped mass or a
one-degree-of-freedom model for thermal analysis, as is the case with DAS and
ORSAT. The improvements introduced by DRAPS with respect to the two com-
puter codes described above are due to new object shapes (15 new types and 51
pre-defined motions) and relevant aerodynamic and aerothermal models, new
break-up criteria, and the use of a simple Monte Carlo method for performing the
uncertainty analysis of re-entry of single objects. In addition, DRAPS uses not only
a criterion based on altitude for predicting break-up, as is the case with ORSAT, but
also a criterion based on temperature.
By contrast, the spacecraft-oriented method consists in modelling the shape and
the structure of a spacecraft in such a way as to make them as near to the reality as
possible. The SCARAB (SpaceCraft Atmospheric Re-entry and Aerothermal
Breakup, Ref. [106]) computer code is based on this method. To this end, the
components of which a spacecraft is made are modelled hierarchically. The ele-
ments which are at the end of this hierarchical structure are simple objects, such as
spheres, cylinders, plates, and boxes. These objects are organised in an order of
parent–children. They are placed at the real locations with respect to the other
objects, in order to make the modelled distribution similar to the real one. For the
purpose of representing the position and the attitude of a re-entering spacecraft, a
six-degree-of-freedom model is used. The forces and torques of aerodynamic nature
are considered over all surface panels in the aerodynamic module. Pressure and
heating on each surface panel are calculated in a way similar to that of the computer
codes described above. For this purpose, it is necessary to know the free-stream
direction and the inclination angle between this direction and the normal to each
surface panel. A two-dimensional model is used for ablation analysis. When the
heat absorbed by a panel exceeds the maximum capacity of that panel, then a
melting condition occurs, and the panel generates fragments. Likewise, when the
mechanical stresses acting upon an object exceed a maximum level, then that object
is broken up. SCARAB has been used for many space programmes, such as the
German satellite ROSAT, the European Ariane 5 launcher, the Italian satellite
BeppoSAX, and the Russian MIR space station.
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Chapter 5
Interplanetary Trajectories
By interplanetary trajectory, we mean the path of a spacecraft which travels from

one to another planet of the Solar System. There have been space probes (Voyager
1, Voyager 2, Pioneer 10, Pioneer 11, and New Horizons) coming from the Earth
and going beyond the orbits of the outer planets of the Solar System. The trajec-
tories of some of these probes are shown in the following figure, which is due to the

https://doi.org/10.1007/978-3-319-62220-0_5
760 5 Interplanetary Trajectories
As is well known, the kinetic energy (½mv2) of a spacecraft of mass m meant to

leave the Solar System must be at least equal to its potential energy (−mlS/r) due to
the attractive force exerted by the Sun, in order for the mechanical energy, E, of that
spacecraft to be equal to or greater than zero. As has been shown in Sect. 4.1, the
escape velocity, ve, is defined as follows
1 l
E ¼ mv2e m S ¼ 0
2 r
At v = ve, a spacecraft moves along a parabola, whose focus is the centre of mass
of the Solar System. At v > ve, a spacecraft moves along a branch of hyperbola,
which has one of its two foci in the same centre of mass. Therefore, for a spacecraft
subject to the gravitational attraction of the Sun, the escape velocity is
12
2lS
ve ¼
r
where lS = 1.327 1020 m3/s2 [2] is the gravitational parameter of the Sun, and
r is the distance of the spacecraft from the centre of mass of the Solar System.
In case of probes meant to leave the Solar System, the escape velocity is usually
reached by means of gravity assist manoeuvres, as the sequel will show.
The principal force acting on a spacecraft travelling from one to another planet
of the Solar System is the gravitational attraction of the Sun. The other forces acting
upon the spacecraft can be considered as perturbations of this principal force. This
holds during the greatest part of the time taken by the space travel, when the
spacecraft is far enough from the planets of, respectively, departure and arrival.
Except for the orbits of Mercury and Pluto, the orbits of the planets belonging to
the Solar System are ellipses of small eccentricity which lie in planes very near to
the plane of the ecliptic. Five of the six orbital elements of these planets are nearly
constant, except for slight perturbations due to the mutual attractions exerted by the
planets. The sixth orbital element indicates the position of each planet along its orbit
and varies with time in accordance with Kepler’s laws. The positions of the planets
at the time of interest can be computed as has been shown in Sects. 3.14 and 3.15.
A practical example of this computation will be given in the next paragraphs.
In case of space vehicles propelled only by chemical rockets, all manoeuvres are
assumed to be impulsive, because the burn times are negligible in comparison with
the total times which are necessary for interplanetary transfers.
5.2 The Hohmann Ellipse Approximation
The computation of an interplanetary trajectory is a complex procedure, which

requires the numerical integration of the equations of motion. These equations must
take account of the forces which perturb the principal attracting force.
5.2 The Hohmann Ellipse Approximation 761
However, in order to perform a preliminary mission analysis or a feasibility

study, it is sufficient to use simple analytical methods. These methods compute the
magnitude and the direction of the impulses which are necessary in order for a
spacecraft to change instantaneously its trajectory from one orbit to another.
The simplest of these methods assumes the planets of the Solar System to
revolve about the Sun in circular and coplanar orbits. This assumption is not far
from the truth, as shown in the following table [3] which indicates the eccentricity
and the orbital inclination in degrees for each planet with respect to the ecliptic
plane.
Mercury Venus Earth Mars Jupiter Saturn Uranus Neptune

Eccentricity 0.205 0.007 0.017 0.094 0.049 0.057 0.046 0.011
Inclination 7.0 3.4 0.0 1.9 1.3 2.5 0.8 1.8
With the exception of Mercury, all planets have orbits of small eccentricity and
small inclination with respect to the ecliptic plane. Hence, according to this
assumption, the trajectory of a spacecraft and the orbits of all the planets are
considered to lie on the ecliptic plane.
As has been shown in Sect. 4.12, a Hohmann semi-ellipse is the most efficient
trajectory for the transfer of a spacecraft from one to another orbit. The following
figure, due to the courtesy of NASA/JPL [4], shows two Hohmann ellipses for
interplanetary transfers from, respectively, the Earth to Mars (left) and the Earth to
Venus (right).
In the Hohmann transfer orbit shown above on the left-hand side, the perihelion
is the departure point, and the aphelion is the arrival point. The radius of perihelion
is the radius of the orbit (assumed circular) of the Earth, and the radius of aphelion
is the radius of the orbit (assumed circular) of Mars.
The orbital velocity, vE, of the Earth with respect to the Sun results from the
vis-viva integral v2/l = 2/r − 1/a, as follows
12
lS
vE ¼
rE
where lS is the gravitational parameter of the Sun, and rE is the radius of the
circular orbit of the Earth.
The velocity, vd, of the spacecraft with respect to the Sun at the departure point
(subscript d) also results from the vis-viva integral
12
2 1
vd ¼ l S
r E aH
where aH = ½(rE + rM) is the major semi-axis of the Hohmann ellipse, and rM is the
radius of the circular orbit of Mars.
This velocity is greater than the orbital velocity of the Earth with respect to the
Sun. Therefore, in order to insert the spacecraft into a Hohmann transfer ellipse at
the departure point, an impulse
Dvd ¼ vd vE
must be given to the spacecraft in the direction of the orbital velocity of the Earth.
At the arrival point (subscript a), the velocity of the spacecraft with respect to the
Sun
12
2 1
va ¼ lS
r M aH
is less than the orbital velocity, vM, of Mars with respect to the Sun
12
lS
vM ¼
rM
Therefore, in order to insert the spacecraft from its transfer orbit into the circular
orbit of Mars, an impulse
Dva ¼ vM va
must be given to the spacecraft in the direction of the orbital velocity of Mars.
The two impulses Dvd and Dva are, both of them, to be given to the spacecraft in
the same direction as that of the orbital velocity of the planet, because the velocity
of Mars at any point of its circular orbit is greater than the velocity of the spacecraft
at the point of arrival (aphelion) of its semi-elliptic transfer orbit.
5.3 The Departure and Arrival Times 763
5.3 The Departure and Arrival Times
A spacecraft to be used for an interplanetary mission cannot leave the orbit of the
planet of departure at any time, because that spacecraft is meant to reach, at the end
of its interplanetary trajectory, not only the orbit of the planet of destination, but the
planet itself. To this end, it is necessary to design a rendezvous between the
spacecraft and the target planet.
In order for a rendezvous to occur at the moment of arrival, the position of the
target planet along its orbit around the Sun, at the moment in which the spacecraft
leaves the planet of departure, must be such that the target planet arrives at the
intersection point between the apsidal line of the transfer semi-ellipse and its own
orbit at the same time as does the spacecraft.
With reference to the following figure, due to the courtesy of NASA/JPL [4], let
us consider a Hohmann ellipse, along which a spacecraft travels from the Earth to
Venus.
Both of these planets are assumed here to revolve about the Sun in circular orbits
which lie in the ecliptic plane.
In this case, any straight line lying in the ecliptic plane and passing through the
centre of mass of the Sun may be chosen as the apsidal line of a Hohmann ellipse.
Let /E0 and /V0 be the true anomalies of, respectively, the Earth and Venus at time
t0. At any other time t, the true anomalies of these planets are, respectively
/E ¼ /E0 þ nE t
/V ¼ /V0 þ nV t
1= 1=
where nE = ðlS =r E 3 Þ 2 and nV = ðlS =r V 3 Þ 2 are the mean motions of, respectively,
the Earth and Venus, rE and rV are the radii of their respective circular orbits, and
lS is the gravitational parameter of the Sun.
As has been shown in Sect. 4.9, the phase angle, w, between the position vectors
of the two planets at any given time is defined as follows
w ¼ /V /E
Hence, there results
w ¼ /V /E ¼ /V0 þ nV t ð/E0 þ nE tÞ ¼ /V0 /E0 þ ðnV nE Þt

¼ w0 þ ðnV nE Þt
where w0 = /V0 − /E0 is the phase angle between the position vectors of the two
planets at time t0. In the case considered above, the orbit of the planet of destination
(Venus) lies inside the orbit of the planet of departure (the Earth). Therefore,
according to the third law of Kepler, nV − nE > 0. In other words, Venus moves
counterclockwise with respect to the Earth.
The opposite is true in case of the planet of destination (e.g. Mars) being outside
the orbit of the planet of departure (e.g. the Earth). In this case, nM − nE < 0,
because Mars moves clockwise with respect to the Earth.
The phase angle between two planets depends on time, as has been shown
above. Since the mean motion of Venus is higher than that of the Earth, then Venus
overtakes the Earth after the lapse of a given time; that is, the phase angle of Venus
with respect to the Earth returns to the value w0 after the lapse of a given time. The
synodic period of Venus, Tsyn, is defined by the following expression
w0 þ 2p ¼ w0 þ ðnV nE ÞTsyn
This expression, solved for Tsyn, yields
2p
Tsyn ¼
nV nE
By contrast, the synodic period of Mars is
2p
Tsyn ¼
nE nM
because the mean motion of the Earth is higher than that of Mars, in accordance
with the third law of Kepler. The synodic period of a planet is the amount of time
required in order for that planet to return to the same point in its orbit with respect to
the Earth and the Sun. Since the orbital period, TP, and the mean motion, nP, of any
given planet P are such that
2p
TP ¼
nP
then the two preceding expressions for Tsyn may also be written as follows
TV TE
Tsyn ¼
TE TV
and
TM TE
Tsyn ¼
TM TE
In both cases, Tsyn indicates the synodic period of the planet of arrival (either
Venus or Mars) with respect to the planet of departure (the Earth).
Let us consider an interplanetary mission in which a spacecraft leaves the Earth
and goes to Mars along a Hohmann semi-ellipse. After the lapse of a certain time
from the time of arrival, the same spacecraft leaves Mars and returns to the Earth
along another Hohmann semi-ellipse. This mission is illustrated in the following
figure, where the two parts, go and back, of the space travel are indicated,
respectively, on the left-hand side and on the right-hand side. The direction of
motion in the two circular orbits is counterclockwise. This also holds with the
Hohmann semi-ellipse, which has the centre of mass of the Sun at one of its foci.
The first Hohmann semi-ellipse is designed in order for the rendezvous between
the spacecraft and Mars to occur at the aphelion. Likewise, the second Hohmann
semi-ellipse is designed in order for the rendezvous between the spacecraft and the
Earth to occur at the perihelion.
The major semi-axes of the two Hohmann semi-ellipses have the same value,
which results from
1
aH ¼ ðrE þ rM Þ
2
where rE and rM are the radii of the orbits of the two planets, the Earth and Mars.
The time necessary to cover one the two parts (go and back) of the total travel is
equal to the time necessary to cover the other part. This time results from
3 12
aH
tEM ¼ p
lS
where lS is the gravitational parameter of the Sun.

In order for a rendezvous between the spacecraft and Mars to occur after the time
tEM, it is necessary that, during the time tEM, Mars arrives, in its motion along its
circular orbit, at the aphelion of the Hohmann semi-ellipse along which the
spacecraft travels.
Since the mean motion of Mars is nM = (lS/r3M)½, then the angle through which
Mars travels, when the spacecraft travels through an angle of p radians, is nMtEM.
With reference to the preceding figure, let w0 be the phase angle between the
Earth and Mars at the time of departure. In order for a rendezvous to occur at the
time of arrival, the following condition must be satisfied
w0 þ nM tEM ¼ p
The preceding equation determines the value of w0.

Let wf be the phase angle between the Earth and Mars at the time of arrival. The
value of wf results from the equation written above, that is, from
w ¼ w0 þ ðnV nE Þt
which relates to a transfer from the Earth to Venus.

In the present case, we consider a transfer from the Earth to Mars, and therefore
wf ¼ w0 þ ðnM nE ÞtEM
Since w0 = p − nMtEM, then
wf ¼ p nM tEM þ ðnM nE ÞtEM ¼ p nE tEM
Since the orbital motion of Mars is slower than that of the Earth, then Mars is
behind the Earth, at the time of arrival, by an angle equal to wf. So much for the go
part of the total travel.
As to the back part, let w*0 be the phase angle between the Earth and Mars at the
time of departure. As has been shown above, the times required to cover the two
parts (go and back) of the total travel have the same value, which results from
3 12
aH
tEM ¼ p
lS
Therefore, the angular distance covered by the Earth during the back part of the
travel is the same as that covered by the Earth during the go part. In other words,
there results
w 0 ¼ wf
which means that the phase angle, w*0, between the Earth and Mars at the time of
departure in the back part of the travel must be equal in magnitude, but opposite in
sign, to the phase angle, wf, between the Earth and Mars at the time of arrival in the
go part of the travel.
Let twait be the so-called wait time, which is the time necessary for the phase
angle between the Earth and Mars to change from wf to −wf. By setting t = 0 for the
time of arrival at Mars, the preceding equation w = w0 + (nM − nE)t becomes
w ¼ wf þ ðnM nE Þt
According to the definition of twait given above, the following equality must hold
wf ¼ wf þ ðnM nE Þtwait
This equality, solved for twait, yields
2wf
twait ¼
nM nE
where wf = p − nEtEM. The value of twait is positive, when the value of wf is also
positive, because nM − nE < 0, and is negative otherwise. When the value of twait,
computed as indicated above, is less than zero, then the favourable opportunity for
beginning the back part of the travel has occurred before the time at which the
spacecraft has arrived at the target planet, and will occur again in the future.
In the general case of a space travel between any two planets, it is necessary to
either add 2Np (where N = 0, 1, 2, …) to −2wf, or subtract 2Np from −2wf, in order
for twait to have a positive value. In particular, when the mean motion, nA, of planet
of arrival is less than the mean motion, nD, of the planet of departure, then
2wf 2Np
twait ¼
nA nD
By contrast, when the mean motion, nA, of planet of arrival is greater than the
mean motion, nD, of the planet of departure, then
2wf þ 2Np
twait ¼
nA nD
where the value of N is chosen in order for the value of twait to be positive and as
small as possible.
For a travel Earth–Mars–Earth, there results nA < nD, because nM < nE. The
radii of the orbits of the two planets and the gravitational parameter of the Sun are
(from Ref. [3]):
rE ¼ 1:496 1011 m
rM ¼ 2:279 1011 m
lS ¼ 1:327 1020 m3 =s2
The major semi-axis of the Hohmann semi-ellipse is
rE þ rM 1:496 1011 þ 2:279 1011

aH ¼ ¼ ¼ 1:888 1011 m
2 2
The travel time to go from the Earth to Mars is
3 12 " #1
3 2
aH ð1:888 1011 Þ
tEM ¼p ¼ 3:1416 ¼ 2:236 107 s
lS 1:327 1020
The same time, expressed in days, is tEM = 258.8 days.

The mean motion of the Earth and the mean motion of Mars are, respectively,
12 " #12

lS 1:327 1020
nE ¼ 3 ¼ ¼ 1:991 107 rad/s
rE ð1:496 1011 Þ3
12 " #12
lS 1:327 1020
nM ¼ 3 ¼ ¼ 1:059 107 rad/s
rM ð2:279 1011 Þ3
The phase angle, w0, between the Earth and Mars when the spacecraft leaves the
Earth depends on the condition of rendezvous between the spacecraft and Mars at
the aphelion of the Hohmann semi-ellipse. As has been shown above, this condition
requires the following value for the phase angle
w0 ¼ p nM tEM ¼ p 1:059 107 2:236 107 ¼ 0:7737 rad ¼ 44 :33

This opportunity occurs once in each synodic period. As has been shown above,
the synodic period of Mars is given by
2p 2 3:1416
Tsyn ¼ ¼ ¼ 6:741 107 s
nE nM 1:991 107 1:059 107
The same time, expressed in days, is Tsyn = 780.2 days.

The phase angle, wf, between the Earth and Mars, when the spacecraft arrives at
Mars, is
wf ¼ p nE tEM ¼ p 1:991 107 2:236 107 ¼ 1:311 rad
Since nM < nE, then the wait time is given by
2wf 2Np
twait ¼
nM nE
The minimum value of N, for which the wait time is greater than zero, is 1.
For N = 1, there results
2 ð1:311Þ 2p
twait ¼ ¼ 3:929 107 s
1:059 107 1:991 107
The same time, expressed in days, is twait = 454.7 days.

Since the travel time necessary to go from the Earth to Mars along a Hohmann
semi-ellipse is tEM = 258.8 days and the wait time is twait = 454.7 days, then the
minimum total time required for this type of go-and-back mission to Mars is
ttotal ¼ 258:8 þ 454:7 þ 258:8 ¼ 972:3 days
5.4 The Spheres of Influence
In the preceding paragraphs, some transfer trajectories have been computed by

means of a two-body model, in which the Sun has been considered the main
attracting body, and the perturbations induced by all planets have been neglected.
A Hohmann transfer trajectory is a semi-ellipse having the centre of mass of the
Sun at one of its foci. In the vicinity of the planet of arrival, an impulse is necessary
in order to insert a spacecraft from the transfer trajectory into a capture orbit around
the planet of arrival. In order to be captured into an orbit around the planet of
arrival, a spacecraft must decrease its velocity with respect to the planet of arrival
by using a retrograde rocket burn or some other means. When that spacecraft is
meant to land on the target planet, it is necessary to reduce further the velocity by
using a retrograde burn. In case of the target planet having an atmosphere, the final
deceleration may also be performed by aerodynamic braking. Other means, such as
parachutes and further retrograde burns, may also be used [4].
The expression “in the vicinity” used above means that we assume a spacecraft
to be subject to the gravitational attraction exerted by only one principal body at a
time, and therefore the equations holding for the two-body problem, described in
Chap. 1, can be used to describe approximately the motion of that spacecraft.
Another expression which has the same meaning is “within the sphere of influ-
ence”. As long as a spacecraft is within the sphere of influence of the Sun or within
the sphere of influence of a given planet, the solar or the planetary attraction is the
principal force acting upon that spacecraft, and any other force is considered as a
perturbation of the principal force.
Let us consider again a spacecraft, which travels from a planet of departure to
another planet of arrival within the Solar System. When that spacecraft leaves the
sphere of influence of the planet of departure, it enters the sphere of influence of the
Sun. Likewise, when that spacecraft leaves the sphere of influence of the Sun, it
enters the sphere of influence of the planet of arrival. The method based on this
approximation is called the patched-conic method, because a spacecraft travelling
within the sphere of influence of a celestial body moves along a conic section
having the centre of mass of that celestial body at one of its foci. The transition
from a geocentric motion to a heliocentric motion is a gradual process, which
occurs not instantaneously but rather within a finite interval of time. During this
time, the gravitational attraction exerted by the Earth on the spacecraft decreases
gradually, and the gravitational attraction due to the Sun increases gradually.
Within the sphere of influence of the Earth, the gravitational attraction due to the
Earth is the principal force acting on a spacecraft, and the gravitational attraction
due to another celestial body is the perturbing force. Therefore, within the sphere of
influence of the Earth, it is advantageous to write the equations of motion of a
spacecraft in a geocentric inertial reference system.
Likewise, within the sphere of influence of the Sun, the gravitational attraction
due to the Sun is the principal force acting on a spacecraft, and the gravitational
attraction due to another celestial body is the perturbing force. Therefore, within the
sphere of influence of the Sun, it is advantageous to write the equations of motion
of a spacecraft in a heliocentric inertial reference system.
The question arises as to which of the two reference systems should be chosen
for any given distance of the spacecraft from the Earth. Laplace proposed to choose
that system of reference for which the ratio of the perturbing force to the principal
force is smaller.
Following Battin [3], with reference to the follow figure, let m1, m2, and m3 be
the masses of, respectively, the Earth, the spacecraft, and the Sun. Let xyz be an
inertial system of reference, whose origin is the centre of mass of the Earth, and
5.4 The Spheres of Influence 771
whose fundamental plane xy is the ecliptic plane. Let XYZ be another inertial system
of reference, whose origin is the centre of mass of the Sun, and whose fundamental
plane XY is the ecliptic plane.
The trajectory of the spacecraft lies in the ecliptic plane. Let r and q be the
position vectors of, respectively, the spacecraft and the Sun in the geocentric inertial
system xyz. Let d = r − q be the position vector of the spacecraft in the heliocentric
inertial system XYZ.
The motion of the spacecraft in the geocentric system of reference is governed

by the following differential equation

r d q
r00 þ Gðm1 þ m2 Þ ¼ Gm 3 þ
r3 d 3 q3
and the motion of the spacecraft in the heliocentric system of reference is governed
by the following differential equation

00 d r q
d þ Gðm2 þ m3 Þ 3 ¼ Gm1 3 3
d r q
where r″ and d″ indicate the second derivatives of, respectively, r and d with respect
to time. The four acceleration vectors, which act upon the spacecraft, are ap 21 , ad 21 ,
ap 23 , and ad 23 . The first two (ap 21 and ad 21 ) of these vectors are considered in the
geocentric system; the other two (ap 23 and ad 23 ) are considered in the heliocentric
system. In this notation, the superscript letter (p or d) indicates the principal
acceleration or the disturbing acceleration, and the subscripts numbers indicate
whether the motion of the spacecraft (subscript 2) is considered in the geocentric
(subscript 1) system or in the heliocentric (subscript 3) system. These four accel-
erations are expressible as follows
m1 þ m2
ap 21 ¼ G ur
r2
1 k
m3 X r
ad 21 ¼G 2 P0k þ 1 ðcos cÞuq P0k ðcos cÞur
q k¼1 q
m2 þ m3
ap 23 ¼ G ud
"d2 #

m1 r 2
ad 23 ¼G 2 uq ur
r q
In the expressions written above, ur, uq, and ud are unit vectors along, respec-
tively, r, q, and d, and P0 k þ 1 (cos c) uq and P0 k (cos c) are the first derivatives of the
Legendre polynomials with respect to their argument (cos c). By c, we mean the
angle between r and q. The Legendre polynomials are used to take account of the
non-spherical shape of the Earth. They have been introduced to the reader in
Sect. 3.2.
The ratio of the disturbing acceleration to the principal acceleration for the mass
m2 relative to the mass m3 is exactly
ad 23 m1 1 2ðr=qÞcos c þ ðr=qÞ2 h 2
i1
4 2
¼ 1 2 ð r=q Þ cos c þ ð r=q Þ
ap 23 m2 þ m3 ðr=qÞ2
The ratio of the disturbing acceleration to the principal acceleration for the mass
m2 relative to the mass m1 is approximately
ad 21 m3 1 h i
¼ ðr=qÞ3 1 þ 3 cos2 c 2 þ O ðr=qÞ4
a 21 m1 þ m2
p
where terms of the order (r/q)4 and higher are neglected.

By equating these two ratios and assuming the value of r/q to be negligible with
respect to unity, there results
1
r m1 ðm1 þ m2 Þ 5 1
¼ 1 þ 3 cos2 c 10
q m3 ðm2 þ m3 Þ
The term (1 + 3 cos2c)1/10 is at most equal to 1.15. In addition, the mass (m2) of
the spacecraft is negligible in comparison with the mass (m1) of the Earth and, a
fortiori, in comparison with the mass (m3) of the Sun. Therefore, the preceding
expression may be approximated as follows
25
m1
r¼q
m3
provided that m1 is much smaller than m3.

5.4 The Spheres of Influence 773
The preceding equation determines the radius of a sphere, whose centre is the
centre of mass of the Earth, and whose surface is the locus of all points, for which
the ratio of the disturbing acceleration to the primary acceleration is the same for
either of the two systems of reference, geocentric and heliocentric, in which the
motion of a spacecraft is considered.
This sphere is called the sphere of influence of the Earth (of mass m1) with
respect to the Sun (of mass m3). The coefficient q is the major semi-axis of the orbit
of the Earth around the Sun. When a spacecraft moves inside this sphere, it is
advantageous to use a geocentric reference system as the origin of the co-ordinates,
whereas, when a spacecraft moves outside the same sphere, it is advantageous to
use a heliocentric reference system.
The same line of reasoning holds when the mass (m1) and the major semi-axis
(q) of any other planet are considered with respect to the mass (m3) of the Sun. The
following table (from Ref. [3]) indicates the radii of the spheres of influence
(SOI) in km for the planets of the Solar System.
Planet Major semi-axis (km) Mass ratio Planet/Sun Radius of SOI (km)
Mercury 57.9 106 0.16595 10−6 112000
Venus 108.2 106 2.4491 10−6 616350
Earth 149.6 106 3.0023 10−6 924510
Mars 227.9 106 0.32286 10−6 577230
Jupiter 778.6 106 954.49 10−6 48219000
Saturn 1433.5 106 285.64 10−6 54793000
Uranus 2872.5 106 43.651 10−6 51791000
Neptune 4495.1 106 51.295 10−6 86451000
When the term (1 + 3 cos2c)1/10 is taken into account, then the surface of
influence of a given planet with respect to the Sun is an oblate spheroid instead of a
sphere.
5.5 The Patched-Conic Approximation
This approximation assumes that a spacecraft, when travelling outside the sphere of
influence of a planet, moves along a conic section having the centre of mass of the
Sun at one of its foci. In other words, the gravitational attraction due to the Sun is
assumed to be the only force acting upon the spacecraft. In addition, outside the
spheres of influence of the planets, the radii of these spheres are neglected, and the
planets are considered as point masses. This assumption holds along the interme-
diate (which is also the greatest) segment of an interplanetary trajectory. Likewise, a
spacecraft travelling inside the sphere of influence of a planet of departure is
assumed to move along a conic section having the centre of mass of that planet at
one of its foci. This holds along the initial segment of an interplanetary trajectory,
before the spacecraft leaves the sphere of influence of the planet of departure. This
also holds along the final segment of the same trajectory, after the spacecraft enters
the sphere of influence of the planet of destination.
In order to use the method described above, we first consider the intermediate
segment of the trajectory. This segment is a branch of ellipse (e.g. a Hohmann
semi-ellipse, as has been shown in the preceding paragraphs), which has the centre
of mass of the Sun at one of its foci, and the appropriate positions of the two planets
along their respective orbits at its endpoints.
The spacecraft moves along this segment from the sphere of influence of the
planet of departure to the sphere of influence of the planet of arrival. When the
spacecraft reaches the sphere of influence of the planet of departure at the time td,
the orbital velocity, vd, of the spacecraft is computed in an inertial heliocentric
system of reference. This is the velocity needed by the spacecraft to be inserted
from the initial segment into the intermediate segment of trajectory. The hyperbolic
excess velocity (see Sect. 1.7) of the spacecraft with respect to the planet of
departure at t = td results from a vector subtraction
vd1 ¼ vd vPd
where vPd is the orbital velocity of the planet of departure with respect to the Sun at
t = td. Likewise, when the spacecraft reaches the sphere of influence of the planet of
arrival at the time ta, the orbital velocity, va, of the spacecraft is computed again in
the same heliocentric system of reference. This is the velocity at which the
spacecraft is inserted from the intermediate segment into the final segment of
trajectory. The hyperbolic excess velocity of the spacecraft with respect to the
planet of arrival at t = ta results from a vector subtraction
va1 ¼ va vPa
where vPa is the orbital velocity of the planet of arrival with respect to the Sun at the
time t = ta. By so doing, we put together three segments of trajectory, each of which
belongs to an individual conic section.
5.6 The Departure of a Spacecraft from a Planet
A spacecraft meant to leave a given planet of departure must reach the surface of
the sphere of influence of that planet at a hyperbolic excess velocity, vd∞, of
magnitude, vd∞, greater than zero. If vd∞ were equal to zero, the spacecraft would
reach the surface of the sphere of influence, but would remain in the orbit of the
planet of departure, and therefore would not be inserted into a transfer ellipse
having the centre of mass of the Sun at one of its foci.
5.6 The Departure of a Spacecraft from a Planet 775
With reference to the following figure, we consider a spacecraft which leaves the
sphere of influence of a planet of departure (e.g. the Earth) and goes towards the
sphere of influence of an outer planet (e.g. Mars).
At the point in which the escape hyperbola intersects the sphere of influence of
the planet of departure, the velocity vector, vd, of the spacecraft with respect to the
Sun is parallel to an asymptote of the escape hyperbola. The direction of this
asymptote, at that point, is the same as that of the orbital velocity vector, vPd, of the
planet of departure with respect to the Sun. This is because the transfer trajectory, in
the further segment going from the sphere of influence of the planet of departure to
the sphere of influence of the planet of arrival, is a Hohmann semi-ellipse.
The magnitude of the velocity vector, vd, of the spacecraft with respect to the
Sun at the moment of departure results from the vis-viva integral
12
2 1
vd ¼ lS
rPd aH
where lS is the gravitational parameter of the Sun, aH = ½(rPd + rPa) is the major
semi-axis of the Hohmann semi-ellipse, and rPd and rPa are the radii of the circular
orbits of, respectively, the planet of departure and the planet of arrival.
The magnitude, vd∞, of the hyperbolic excess velocity of the spacecraft with
respect to the planet of departure results from the vis-viva integral
v2 2 1
¼
lPd r a
where lPd is the gravitational parameter of the planet of departure, r is the distance
of the spacecraft from the centre of mass of the planet of departure, and a is the
transverse semi-axis of the escape hyperbola. By setting r equal to infinity, vd∞
results from
12
lPd
vd1 ¼
ðaÞ
It is to be remembered (see Sect. 1.2) that the semi-transverse axis, a, of a

hyperbola has a negative value, so that the value of −a is positive.
On the other hand, the semi-latus rectum, p, of the escape hyperbola is
h2
p ¼ a 1 e2 ¼
lPd
where e > 1 is the eccentricity of the escape hyperbola, and h is the moment of
momentum per unit mass of the spacecraft.
The preceding equation, solved for lPd/(−a), yields
lPd l2 2
¼ Pd e 1
ðaÞ h 2
hence
1
lPd 2 lPd 2 1
vd1 ¼ ¼ e 1 2
ðaÞ h
Let rP be the periapsis radius in the escape hyperbola. Remembering the general
equation of a conic section
p
r¼
1 þ e cos /
and setting / = 0 (i.e. cos / = 1), the radius of periapsis, rP, results from
p
rP ¼
1þe
By substituting p = rP(1 + e) into
v2d1 1 e2 1
¼ ¼
lPd ðaÞ p
there results
rP v2d1
e¼ þ1
lPd
which is the value of the eccentricity of the escape hyperbola.

This value, in turn, substituted into the preceding equation
lPd 2 1
vd1 ¼ e 1 2
h
makes it possible to compute the value of h as follows

12
l
h ¼ rP v2d1 þ 2 Pd
rP
Instead of h, it is possible to compute the value of p = h2/lPd as follows
p ¼ rP ðe þ 1Þ
The hyperbolic excess velocity, vd∞, of a spacecraft with respect to the planet of
departure is a constraint imposed by the mission to be performed, because it is the
impulse (Dvd) necessary in order for the spacecraft to reach the required velocity
[lS(2/rPd − 1/aH)]½ with respect to the Sun, as follows
12 12
2 1 l
vd1 ¼ Dvd ¼ lS S
rPd aH rPd
where lS is the gravitational parameter of the Sun, aH = ½(rPd + rPa) is the major
semi-axis of the Hohmann semi-ellipse, and rPd and rPa are the radii of the circular
orbits of, respectively, the planet of departure and the planet of arrival.
Therefore, when the periapsis radius, rP, has been chosen by the designer, then
the eccentricity, e, and the semi-latus rectum, p, of the escape hyperbola are also
determined, as has been shown above.
The velocity, vP, of the spacecraft at periapsis results from the vis-viva integral
v2 2 1
¼
lPd r a
By substituting r = rP, v = vP and −1/a = (e2 − 1)/p in the equation written

v2P 2 e2 1
¼ þ
lPd rP p
Now, by substituting e2 − 1 = vd1 2 ðr P 2 vd1 2 þ 2Þ=lPd and p ¼ rP ðrP vd1 2 =

lPd þ 2Þ into the preceding equation, there results
12
2lPd
vP ¼ þ vd1 2
rP
Before receiving the impulse necessary for departure, the spacecraft revolves
around the planet of departure in a circular parking orbit of radius rP. The velocity
of the spacecraft in this circular orbit is
12
lPd
vc ¼
rP
Therefore, the impulse needed by the spacecraft to be inserted from the parking
orbit into the escape hyperbola is
Dv ¼ vP vc
With reference to the following figure, the angle, a, between the transverse axis
FF′ and the asymptote OE of the escape hyperbola results from
a 1
cos a ¼ ¼
ea e
Since the value of e ¼ rP vd1 2 =lPd þ 1 has been found above, then
0 1

1 B 1 C
a ¼ arccos ¼ arccosB@ 2
C
A
e rP vd1
þ1
lPd
This is also the angle between the transverse axis of the escape hyperbola and the
orbital velocity vector, vPd, of the departure planet with respect to the Sun.
The plane on which the escape hyperbola lies is any plane containing the centre
of mass of the planet of departure and the hyperbolic excess velocity vector, vd∞.
Curtis [5] points out that the rotation of this hyperbola about an axis which passes
through the centre of mass of the planet of departure and is parallel to vd∞ deter-
mines a surface of revolution, which is the locus of all possible departure hyper-
bolas. Likewise, the rotation of the periapsis, P, of the departure hyperbola about
the axis indicated above determines a circumference of radius rP, which is the locus
of all possible injection points from a circular parking orbit around the planet of
departure into a hyperbola directed towards the planet of arrival.
On the other hand, the plane on which this parking orbit (or the trajectory of direct
ascent) lies must contain the axis indicated above and the launch site at the time of
launch. Let u* be the geocentric latitude of the launch site on the surface of the Earth.
In case of a prograde orbit, the possible inclination angle of the orbit with respect to
the equator ranges from a minimum value of u* (in case of a launch due East) to a
maximum value of 90°. Of course, this holds in case of no safety constraints limiting
the range indicated above. For example, as has been shown in Sect. 4.4, the safety
constraints relating to Cape Canaveral (u* = 28°.5) are Amax 35° and Amin 120°.
Therefore, in case of a prograde orbit, these constraints limit the possible
inclinations for the departure hyperbola to a range between a minimum value of
28°.5 (in case of a launch due East) and a maximum value of 55°.0, as shown by the
following figure, which is due to the courtesy of Space Florida [6].
Since the obliquity of the ecliptic is about 23°.5, then the maximum value for a
is 55°.0 − 23°.5 = 31°.5.
Therefore, when the departure of a spacecraft from the Earth to an outer planet
(such as Mars) occurs on the dark side of the Earth, then the maximum value for a
is 31°.5 counterclockwise, and when the same departure occurs on the side of the
Earth illuminated by the Sun, then the maximum value for a is 31°.5 clockwise, as
In both cases, a is the angle between the transverse axis and the asymptote of the
departure hyperbola (or, which is the same, the angle between the transverse axis of
the departure hyperbola and the direction of the velocity vector, vE, of the Earth
with respect to the Sun at the moment of departure).
Let us consider again the case of a spacecraft launched from Cape Canaveral into
a circular parking orbit around the Earth. Since the value of a varies between about
31°.5 counterclockwise and about 31°.5 clockwise, then the launch window for
each day is limited by the two times in which the launch site, due to the rotation of
the Earth about its axis, passes through the two positions indicated above. After the
spacecraft has been inserted into a parking orbit, a burn to reach the departure
hyperbola can occur at the completion of each orbit.
The case considered above concerns the departure of a spacecraft from the Earth
to an outer planet (such as Mars). By contrast, a spacecraft going from the Earth to
an inner planet (such as Venus) must have a velocity, vd, with respect to the Sun
smaller in magnitude than the orbital velocity, vPd, of the Earth.
To this end, a spacecraft going from the Earth to Venus must come out from the
sphere of influence of the Earth on the side opposite to that of a spacecraft travelling
from the Earth to Mars, and must also have a hyperbolic excess velocity vector,
vd∞, directed oppositely to vPd, as shown in the following figure.
An example given below illustrates the trajectory of a spacecraft leaving the Earth
from a circular parking orbit of radius rP = 6378 + 300 km and going to Mars.
Let lS = 1.327 1011 km3/s2 and lE = 3.986 105 km3/s2 [2] be the gravi-
tational parameters of, respectively, the Sun and the Earth. Using the same source of
information, let rM = 227.9 106 km and rE = 149.6 106 km be the radii of the
orbits of, respectively, Mars and the Earth.
As has been shown above, the hyperbolic excess velocity of the spacecraft with
respect to the Earth results from
12 12
2 1 l
vE1 ¼ Dvd ¼ lS S
r E aH rE
where aH = ½(rE + rM) is the major semi-axis of the Hohmann semi-ellipse.

The values of aH and vE∞ are
rE þ rM 149:6 106 þ 227:9 106

aH ¼ ¼ ¼ 188:75 106 km
2 2
12 12 12
2 1 l 2 1
vE1 ¼ lS S ¼ 1:327 1011
r E aH rE 149:6 106 188:75 106
1
1:327 1011 2
¼ 2:943 km/s
149:6 106
The spacecraft, when revolving about the Earth in a circular parking orbit of
radius rP, has a velocity of
12 1
lE 3:986 105 2
vc ¼ ¼ ¼ 7:726 km/s
rP 6678
The velocity of the spacecraft at the perigee of the escape hyperbola is

12 12
2lE 2 3:986 105
vP ¼ þ vE1 ¼
2
þ 2943 ¼ 11:315 km/s
2
rP 6678
The impulse needed by the spacecraft to be inserted into an escape hyperbola

directed towards Mars results from subtracting the velocity vc from vP, as follows
Dv ¼ vP vc ¼ 11:315 7:726 ¼ 3:590 km=s
The angle between the transverse axis of the departure hyperbola and the
direction of the velocity vector of the Earth at the moment of departure is
0 1 0 1
B 1 C B 1 C
a ¼ arccosB
@
C ¼ arccosB
A @6678 2:9432
C ¼ 29 :16
A
rP v2E1
þ1 þ1
lE 3:986 105
As has been shown above, the perigee of the departure hyperbola may be placed
either on the illuminated side or on the dark side of the Earth. In case of a prograde
orbit, the perigee is placed on the dark side.
As has been shown in Sect. 4.19, the mass, Dm, of propellant necessary for
inserting the spacecraft into the departure hyperbola may be computed by means of
Tsiolkovsky’s equation

Dv
Dm ¼ m1 1 exp
g0 Isp
provided that the total mass, m1, of the spacecraft and the specific impulse, Isp, of
the rocket motor are known. Assuming Isp = 300 s, there results

Dm Dv 3:590
¼ 1 exp ¼ 1 exp 0:7
m1 g0 Isp 0:00981 300
In other words, before the burn, about 700 kg of propellant are necessary out of
1000 kg of total mass of the spacecraft.
The following section of this paragraph evaluates the effects, which small errors
in either position or velocity or both of them, committed at the moment of burn,
may have on the trajectory. To this end, we consider a transfer from the Earth to an
outer planet, such as Mars. With reference to the following figure, let rE and rM be
the radii of the circular orbits of, respectively, the Earth and Mars.
Let aH = ½(rE + rM) be the major semi-axis of the Hohmann semi-ellipse having
the centre of mass of the Sun at one of its foci.
Since the radius of aphelion of this semi-ellipse is equal to the orbital radius of
Mars, then substituting r = rM and / = p into the general equation of a conic section
r = p/(1 + e cos/) yields
p
rM ¼
1e
where p = h2/lS and e = (rM − rE)/(rM + rE) are, respectively, the semi-latus rectum
and the eccentricity of the Hohmann semi-ellipse, h is the moment of momentum
per unit mass of the spacecraft, and lS is the gravitational parameter of the Sun.
The value of h is constant along the Hohmann semi-ellipse and results from
h ¼ r E vd
where vd = [lS(2/rE − 1/aH)]½ is the velocity of the spacecraft at the point of

departure (subscript d) with respect to the Sun.
By substituting these expressions into rM = p/(1 − e) and solving for rM, there
results
ðrE vd Þ2
rM ¼
2lS rE v2d
The expression written above shows how the real radius of apogee (rM) reached
by the spacecraft at the end of the Hohmann semi-ellipse depends on its real
velocity (vd) with respect to the Sun at the point of burnout. Should the real value of
vd be different from the correct value [lS(2/rE − 1/aH)]½, then the real value of rM
would also differ from the correct value 227.9 106 km.
Let drM be the change in rM due to a small variation dvd of vd. This change in rM
results from

drM
drM ¼ dvd
dvd
Taking account of the expression rM = (rEvd)2/(2lS − rEv2d), there results

" #
4rE2 lS vd
drM ¼ dvd
ð2lS rE v2d Þ2
Dividing the preceding expression by rM = (rEvd)2/(2lS − rEv2d), there results

drM 4lS dvd
¼
rM 2lS rE v2d vd
The velocity vd = [lS(2/rE − 1/aH)]½ of the spacecraft at the point of departure

with respect to the Sun is the sum of the orbital velocity vE of the Earth with respect
to the Sun and the hyperbolic excess velocity vE∞ of the spacecraft with respect to
the Earth, that is,
vd ¼ vE þ vE1
Remembering the expression written above,

12
2lE
vP ¼ þ vE1
2
rP
where lE is the gravitational parameter of the Earth, and vP is the velocity of the
spacecraft at the perigee of the departure hyperbola, and solving for vE∞, there
results
1
2l 2
vE1 ¼ v2P E
rP
and hence
1
2lE 2
vd ¼ vE þ vE1 ¼ vE þ vP
2
rP
The change in vd due to variations of the perigee radius rP and of the perigee
velocity vP results from

@vd @vd
dvd ¼ drP þ dvP
@rP @vP
By executing the partial derivatives, there results
@vd lE lE
¼ 12 ¼ r 2 v
@rP 2l P E1
rP2 v2P E
rP
@vd vP vP
¼ 1 ¼
@vP 2l 2 vE1
v2P E
rP
Hence

lE vP
dvd ¼ drP þ dvP
rP2 vE1 vE1
Remembering again the expression

12
2lE
vP ¼ þ v2E1
rP
there results
2lE 2lE
þ v2E1 þ vE1
vP v2P r r v
¼ ¼ P ¼ P E1
vE1 vP vE1 vP vE1 vP
and therefore
0 1
2lE
þ v
lE BrP vE1 E1 C
dvd ¼ drP þ B
@
CdvP
A
rP2 vE1 vP

dvd lE drP 2lE vE1 dvP
¼ þ þ
vd rP vd vE1 rP rP vd vE1 vd vP

drM 4lS dvd
¼
rM 2lS rE v2d vd
yields

drM 4lS lE drP 2lE vE1 dvP
¼ þ þ
rM 2lS rE v2d rP vd vE1 rP rP vd vE1 vd vP
As an example of application of the equation indicated above, let us consider

again a spacecraft starting from a circular orbit of radius rP = 6378000 + 300000 m
around the Earth and directed to Mars. The radii of the orbits of the two planets and
the gravitational parameter of the Sun are (from Ref. [3]):
rE ¼ 1:496 1011 m
rM ¼ 2:279 1011 m
lS ¼ 1:327 1020 m3 =s2
The major semi-axis of the Hohmann semi-ellipse is
rE þ rM 1:496 1011 þ 2:279 1011

aH ¼ ¼ ¼ 1:8875 1011 m
2 2
The velocity (in m/s) of the spacecraft with respect to the Sun at the moment of
departure is
12 12
2 1 2 1
vd ¼ lS ¼ 1:327 10
20

rE aH 1:496 1011 1:8875 1011
¼ 32726
The hyperbolic excess velocity of the spacecraft with respect to the Earth is
12 1
lS 1:327 1020 2
vE1 ¼ vd ¼ 32726 ¼ 2942:9 m/s
rE 1:496 1011
The velocity of the spacecraft at the perigee of the escape hyperbola is

12 12
2lE 2 3:986 1014
vP ¼ þ vE1 ¼
2
þ 2942:9 ¼ 11315 m/s
2
rP 6678000
By substituting these values into the following equation

drM 4lS lE drP 2lE vE1 dvP
¼ þ þ
rM 2lS rE v2d rP vd vE1 rP rP vd vE1 vd vP
there results

drM drP dvP
¼ 3:127 þ 6:708
rM rP vP
In other words, a variation of 0.01% in the velocity, vP, at the moment of burnout
changes the radius of the target by 0.0001 6.708 = 0.0006708 or 0.06708%.
Since rM = 2.279 1011 m, then 0.0006708 2.279 1011 = 1.529 108 m =
152900 km. A variation of 0.01% in the radius, rP, at the moment of burnout
changes the radius of the target by 0.0001 3.127 = 0.0003127 or 0.03127%. The
corresponding error in the target radius is 0.0003127 2.279 1011 =
7.126 107 m = 71260 km.
For the reason shown above, the small errors which may be committed in the
early phases of a space mission must be corrected in the course of the semi-elliptic
transfer trajectory between the two planets.
5.7 The Arrival of a Spacecraft at a Planet
The velocity vector of a spacecraft with respect to a planet of arrival, at the moment
of entering the sphere of influence of that planet, is the hyperbolic excess velocity,
va∞. In case of a transfer from the Earth to an outer planet (e.g. to Mars), the
velocity vector, va, of the spacecraft with respect to the Sun is less in magnitude
than the velocity vector, vPa, of that planet with respect to the Sun. This case is
Therefore, a spacecraft directed from the Earth to an outer planet enters the
sphere of influence of the planet of arrival on the same side as that of the Sun.
In case of a transfer trajectory being a Hohmann semi-ellipse, the two vectors
indicated above are parallel, and therefore, the magnitude of the hyperbolic excess
velocity vector results from
va1 ¼ vPa va
By contrast, in case of a transfer from the Earth to an inner planet (e.g. to

Venus), the velocity vector, va, of the spacecraft with respect to the Sun is greater in
magnitude than the velocity vector, vPa, of that planet with respect to the Sun. This
case is shown in the following figure.
5.7 The Arrival of a Spacecraft at a Planet 789
Therefore, a spacecraft directed from the Earth to an inner planet enters the
sphere of influence of the planet of arrival on the side opposite to that of the Sun.
In case of a transfer trajectory being a Hohmann semi-ellipse, there results
va1 ¼ va vPa
In both cases, the trajectory followed by the spacecraft after entering the sphere
of influence of the planet of arrival depends on the type of mission planned by the
designer. When the spacecraft is meant to either hit the surface or enter the
atmosphere of the planet, then the distance, d, between the asymptote of arrival and
the occupied focus of the hyperbola must be such that the radius of periapsis, rP, of
the hyperbola is nearly equal to the radius of the planet. When the spacecraft is
meant to change its trajectory from a hyperbola to an ellipse having one of its foci
in the centre of mass of the planet, then the distance indicated above must be chosen
so that the burn relating to the capture manoeuvre occurs when the spacecraft is at
the proper distance from the surface of the planet.
When the spacecraft is meant to continue along its hyperbolic trajectory after
reaching the point, P, of periapsis, in such a way as to come out from the sphere of
influence of the planet of arrival, then the spacecraft leaves the sphere of influence
at a hyperbolic excess velocity, va∞, which is equal in magnitude to the hyperbolic
excess velocity in the point of entrance, and is directed at an angle, d, with respect
to the asymptote of entrance, as shown in the preceding figures.
This angle is called the flyby turn-angle and depends on the eccentricity, e, of the
hyperbola of arrival according to the following definition given in Sect. 1.2:

d 1
sin ¼
2 e
When the magnitude, va∞, of the hyperbolic excess velocity vector and the
radius of periapsis, rP, of the hyperbola of arrival are determined, then the eccen-
tricity of this hyperbola is also determined. It has been found in Sect. 5.6 to be
rP v2a1
e¼ þ1
lPa
where lPa is the gravitational parameter of the planet of arrival.

Combining the two equations written above makes it possible to compute the
flyby turn-angle, as follows
0 1

1 B 1 C
d ¼ 2 arcsin ¼ 2 arcsinB
@ 2
C
A
e rP va1
þ1
lPa
The distance, d, indicated above has been defined as the length of the segment
HF′ whose direction is perpendicular to the asymptote OE of the hyperbola, as
By inspection of this figure, the segments HF′ and EA are equal in length. The
length of EA is the conjugate semi-axis, b, of the hyperbola shown above, whose
equation is
x2 y2
¼1
a2 b2
in a system of Cartesian axes x and y whose origin is the centre O of the hyperbola
and whose plane is the plane of motion, which contains the centre of mass of the
planet of arrival and the hyperbolic excess velocity vector va∞. Since
1
d b ¼ a e2 1 2
1 1 1
p ¼ a e2 1 ¼ a e2 1 2 e2 1 2 ¼ d e2 1 2
then
p
d¼ 1
ð e2 1Þ 2
where p is the semi-latus rectum of the arrival hyperbola, which depends (as has
been shown in Sect. 5.6) on rP and e as follows
p ¼ rP ð1 þ eÞ
Therefore
12
p eþ1
d¼ ¼ rP
ð e 2 1Þ 2
1
e1
Substituting e ¼ rP va1 2 =lPa þ 1 into the preceding equation yields

12
2lPa
d ¼ rP 1þ
rP v2a1
As it is the case with the departure hyperbola described in Sect. 5.6, likewise the
rotation of the hyperbola of arrival about an axis, which passes through the centre of
mass of the planet of arrival and is parallel to the hyperbolic excess velocity vector
va∞, determines a surface of revolution, which is the locus of all possible hyper-
bolas of arrival. In addition, the rotation of the periapsis, P, of the hyperbola of
arrival about the axis indicated above determines a circumference of radius rP,
which is the locus of all possible injection points from the hyperbola of arrival into
a capture orbit around the planet of arrival.
When the purpose of a space mission is the insertion of a spacecraft into an
elliptic orbit around the planet of arrival, then a burn is necessary at the periapsis P,
which is also the vertex of the hyperbola of arrival, in order to change the trajectory
from this hyperbola to an ellipse having the desired eccentricity.
As has been shown in Sect. 5.6 for the hyperbola of departure, likewise the
velocity, vP, of the spacecraft at the periapsis of the hyperbola of arrival is
12
2lPa
vP ¼ þ v2a1
rP
where lPa is the gravitational parameter of the planet of arrival, rP is the radius of
periapsis of the hyperbola of arrival, and va∞ is the magnitude of the hyperbolic
excess velocity vector.
Let e < 1 be the desired eccentricity of the capture ellipse around the planet of
arrival. Let
1 e2
rP ¼ a ¼ að 1 e Þ
1þe
be the radius of periapsis of this capture ellipse, whose major semi-axis is a. Since
the value of rP does not change in the passage from the hyperbola of arrival to the
ellipse of capture, then we substitute a = rP/(1 − e) into the expression of the
vis-viva integral
1
2 1 2
v ¼ lPa
rP a
This yields
1
1þe 2
v ¼ lPa
rP
where v is the velocity of the spacecraft at the periapsis of the capture ellipse, whose
eccentricity has the desired value e < 1.
The impulse needed to change the value of the velocity at the point of periapsis
from vP to v is
12 1
2lPa 1þe 2
Dv ¼ vP v ¼ þ v2a1 lPa
rP rP
which shows that this impulse has its maximum value for a circular orbit of radius
rP around the planet of arrival, and decreasing values for increasing values of the
eccentricity, e, of the capture ellipse.
In order to reduce the expenditure of propulsive energy to the minimum possible
amount, it is necessary to determine the value of rP which corresponds to the
minimum value of Dv. For this purpose, Curtis [5] divides the terms on both sides
of the preceding equation by va∞, as follows
12 1
Dv 2lPa 1þe 2
¼ þ1 lPa 2
va1 rP v2a1 rP va1
The preceding equation, after setting for convenience x = rP va1 2 =lPa , may be
written as follows
12 1
Dv 2 1þe 2
¼ þ1
va1 x x
The value of e is chosen by the designer. Therefore, the function Dv/va∞ depends
only on x. Let (Dv/va∞)′ and (Dv/va∞)″ be, respectively, the first derivative and the
second derivative of the function Dv/va∞ with respect to x. These derivatives are

0 1 12
ðDv=va1 Þ ¼ ð2 þ xÞ þ ð1 þ eÞ x2
1 3
2
2

00 32 3
ðDv=va1 Þ ¼ ð2x þ 3Þð2 þ xÞ ð1 þ eÞ x2
1 5
2
4
Since x = rP va1 2 =lPa cannot be equal to zero, then the value of x, for which the
first derivative vanishes, is the solution of the following equation
1
ð 2 þ xÞ 2 þ
1 1
ð1 þ eÞ2 ¼ 0
2
This value is
1e
x¼2
1þe
This value, substituted into the expression of the second derivative, yields
1
22 ð 1 þ e Þ 3
ðDv=va1 Þ00 ¼ 3
64ð1 eÞ2
The preceding expression is greater than zero when 0 e < 1, that is, when the
trajectory of the desired eccentricity around the planet of arrival is an ellipse.
Summarising, when
rP v2a1 1e
x¼ ¼2
lPa 1þe
for a given value of eccentricity 0 e < 1 chosen by the designer, then Dv has the
minimum value. The radius of periapsis, which corresponds to the minimum value
of Dv, is
2lPa ð1 eÞ
rP ¼
v2a1 ð1 þ eÞ
The preceding figure shows that, for any elliptic orbit, there results
rP 1 e
¼
rA 1 þ e
Therefore, the radius relating to the maximum distance between the spacecraft
and the occupied focus, F, of the desired elliptic orbit requiring the minimum value
of Dv is
1þe l
rA ¼ rP ¼ 2 2Pa
1e va1
In other words, rA is independent of the value chosen for the eccentricity of the
capture orbit around the planet of arrival and is equal to the radius of the circular
orbit corresponding to the minimum value of Dv.
In order to determine the minimum value of Dv, the value computed above
rP v2a1 1e
x¼ ¼2
lPa 1þe
is substituted into the equation

12 1
Dv 2 1þe 2
¼ þ1
va1 x x
This yields
12
1
Dv ¼ va1 ð 1 eÞ
2
The conjugate semi-axis, d, of the hyperbola of arrival has been shown to be

1
2lPa 2
d ¼ rP 1 þ
rP v2a1
Substituting x = rP va1 2 =lPa = 2(1 − e)/(1 + e) into the preceding equation yields
12
2
d ¼ rP
1e
which is the value of the conjugate semi-axis, d, of the hyperbola of arrival cor-
responding to the minimum value of Dv for a capture ellipse having a given
eccentricity e. Of course, Dv and rP have small values for capture ellipses of high
eccentricity (e 1).
As an application of these concepts to a practical case, with reference to the

preceding figure, it is required to determine the minimum impulse necessary in
order to insert a spacecraft, which departs from the Earth and follows a Hohmann
semi-ellipse, into an elliptic orbit of eccentricity e = 0.68 around Venus. It is also
required to determine the radius, rP, of periapsis, the conjugate semi-axis, d, of the
arrival hyperbola, and the angle, a, between the transverse axis and the asymptote
of the arrival hyperbola. The radii of the orbits of the two planets and the gravi-
tational parameter of the Sun are (from Ref. [3]):
rE ¼ 1:496 1011 m
rV ¼ 1:082 1011 m
lS ¼ 1:327 1020 m3 =s2
As has been shown above, the hyperbolic excess velocity vector of the space-
craft with respect to the inner planet of arrival (Venus) has the following magnitude
12 12
2 1 l
vV1 ¼ lS S
r V aH rV
where aH = ½(rE + rV) is the major semi-axis of the Hohmann semi-ellipse and has
the same direction as the orbital velocity of Venus. Therefore, aH and vV∞ are
rE þ rV 1:496 1011 þ 1:082 1011

aH ¼ ¼ ¼ 1:289 1011 m
2 2
12 12 12
2 1 l 2 1
vV1 ¼ lS S ¼ 1:327 1020
r V aH rV 1:082 1011 1:289 1011
1
1:327 1020 2
¼ 2707 m/s ¼ 2:707 km/s
1:082 1011
Remembering that the radius of periapsis of the desired capture orbit having
eccentricity e and requiring the minimum value of Dv is
2lV ð1 eÞ
rP ¼
v2V1 ð1 þ eÞ
there results
2lV ð1 eÞ 2 0:815 3:986 1014 ð1 0:68Þ

rP ¼ ¼ ¼ 1:688 107 m
v2V1 ð1 þ eÞ 27072 ð1 þ 0:68Þ
where lV = 0.815 3.986 1014 m3/s2 is the gravitational parameter of Venus.

This is because the Venus-to-Earth mass ratio is 0.815 [7].
The minimum value of Dv, corresponding to e = 0.68, results from

12 12
1 1
Dv ¼ va1 ð1 eÞ ¼ 2707 ð1 0:68Þ ¼ 1083 m/s
2 2
The conjugate semi-axis, d, of the hyperbola of arrival, corresponding to the

minimum value of Dv for a capture ellipse of eccentricity e = 0.68, results from
12 12
2 2
d ¼ rP ¼ 1:688 107 ¼ 4:221 107 m
1e 1 0:68
The angle between the transverse axis and the asymptote of the arrival hyperbola
results from
0 1

1 B 1 C
a ¼ arccos ¼ arccosB @ 2
C
A
ehyp rP vV1
þ1
lV
where ehyp is the eccentricity of the hyperbola of arrival. In the present case,
0 1 0 1
B 1 C B 1 C
a ¼ arccos@ A ¼ arccosB
@
C ¼ 43 :6
A
rP v2V1 1:688 10 2707
7 2
þ1 þ 1
lV
0:815 3:986 1014
The major semi-axis and the period of the desired elliptic orbit around Venus
are, respectively,
rP 1:688 107
a¼ ¼ ¼ 5:276 107 m
1e 1 0:68
3 12 " 3
#12
a ð5:276 107 Þ
T ¼ 2p ¼ 2 3:1416 ¼ 133610 s ¼ 37:11 h
lV 0:815 3:986 1014
5.8 The Flight of a Spacecraft Past a Planet
A spacecraft directed towards a planet continues, unless hitting the surface or

entering the atmosphere of the planet, along its hyperbolic trajectory with respect to
the planet. The spacecraft reaches the point of periapsis, then goes beyond that
point, and finally leaves the sphere of influence of the target planet. The following
figure shows the trajectory, the asymptotes, and the transverse axis of the hyperbola
along which the spacecraft moves with respect to the planet.
The trajectory shown in the preceding figure is called a leading-side flyby,

because the periapsis, P, of the hyperbola is on the same side as that of the motion
of the planet with respect to the Sun, as shown by the direction of the orbital
velocity vector, v. At the point, Q1, of entrance of the spacecraft into the sphere of
influence of the planet, the velocity, v1, of the spacecraft with respect to the Sun is
equal to the sum of the orbital velocity, v, of the planet with respect to the Sun and
the hyperbolic excess velocity, v1∞, of the spacecraft with respect to the planet, that
is,
v1 ¼ v þ v11
5.8 The Flight of a Spacecraft Past a Planet 799
Likewise, at the point, Q2, of exit of the spacecraft from the sphere of influence
of the planet, the velocity, v2, of the spacecraft with respect to the Sun is equal to
the sum of the orbital velocity, v, of the planet with respect to the Sun and the
hyperbolic excess velocity, v2∞, of the spacecraft with respect to the planet, as
follows
v2 ¼ v þ v21
The variation of velocity of the spacecraft with respect to the Sun is
v2 v1 ¼ ðv þ v21 Þ ðv þ v11 Þ ¼ v21 v11
The hyperbolic excess velocity vectors, v1∞ and v2∞, are directed along the
asymptotes of the hyperbola. Therefore, these vectors form the same angle, a, with
the transverse axis of the hyperbola. They differ one from the other, because one of
them (v1∞) points towards the centre, O, of the hyperbola, whereas the other (v2∞)
points away from the centre. The magnitude of the two vectors is the same (v∞1 =
v∞2 = v∞).
The angle of rotation of v2∞ with respect to v1∞ is the flyby turn-angle d. Hence,
the two difference vectors, v2 − v1 and v2∞ − v1∞, are directed along the transverse
axis of the hyperbola and point away from the periapsis.
The preceding figure, which illustrates a leading-side flyby, shows that the
vector v2 − v1 (or the vector v2∞ − v1∞) has a negative component along the
direction of the velocity vector v of the planet with respect to the Sun.
By contrast, the following figure illustrates a trailing-side flyby, which is
so-called because the periapsis, P, of the hyperbola is on the side opposite to that of
the motion of the planet with respect to the Sun. This figure shows that the vector
v2 − v1 (or the vector v2∞ − v1∞) has a positive component along the direction of
the velocity vector v of the planet with respect to the Sun.
In other words, a leading-side flyby decreases the velocity of a spacecraft with
respect to the Sun, whereas a trailing-side flyby increases that velocity.
With reference to the two preceding figures, let uP and uS be two orthogonal unit
vectors whose directions are, respectively, the orbital velocity, v, of the target planet
with respect to the Sun (i.e. uP = v/v) and the planet-to-Sun direction.
In the point, Q1, where the incoming spacecraft enters the sphere of influence of
the target planet, the velocity, v1, of the spacecraft with respect to the Sun has the
following two components along the directions indicated above
v1 ¼ ðv1 cos h1 Þ uP þ ðv1 sin h1 Þ uS

where v1 is the magnitude of v1, and h1 is the angle, positive counterclockwise,

between v1 and v.
In other words, the magnitude of h1 is the flight path angle c of the spacecraft in
its semi-elliptic trajectory about the Sun at the point Q1 of entrance into the sphere
of influence of the target planet. In addition, v1 cos h1 is the transverse component
of v1, and −v1 sin h1 is the radial component of v1, as follows
h1 lS
v1 cos h1 ¼ r1 /0 1 ¼ ¼ ð1 þ e1 cos /1 Þ
r 1 h1
lS
v1 sin h1 ¼ r 0 1 ¼ e1 sin /1
h1
where e1, /1, and h1 are, respectively, the eccentricity, the true anomaly, and the
moment of momentum per unit mass of the spacecraft in its the semi-elliptic orbit
around the Sun, and lS is the gravitational parameter of the Sun.
The velocity vector of the target planet in its circular orbit about the Sun is
v ¼ v uP
The magnitude of v is v = (lS/rP)½, where rP is the radius of the circular orbit of

the target planet about the Sun.
At the point Q1 of entrance into the sphere of influence of the target planet, the
hyperbolic excess velocity vector of the spacecraft is
v11 ¼ v1 v ¼ ðv11 ÞP uP þ ðv11 ÞS uS
The components of v1∞ along the directions uP and uS are, respectively,
ðv11 ÞP ¼ v1 cos h1 v
ðv11 ÞS ¼ v1 sin h1
The magnitude of v1∞ is

h i12 1
v11 ¼ ðv11 Þ2P þ ðv11 Þ2S ¼ v21 þ v2 2v1 v cos h1 2
Since v1∞ is known, then specifying a value for the radius, rP, of periapsis makes
it possible to determine the eccentricity, e, and the semi-latus rectum, p, of the
hyperbolic trajectory of the spacecraft with respect to the target planet, by using the
expressions shown in Sect. 5.6, that is,
rP v2a1
e¼ þ1
lP
p ¼ r P ð e þ 1Þ
where lP is the gravitational parameter of the target planet.

With reference to the preceding figures, let w1 and w2 be the angles formed by,
respectively, v1∞ at Q1 and v2∞ at Q2 with respect to the velocity vector v of the
planet in its motion about the Sun. According to this definition, the value of w1 at
Q1 results from

ðv11 ÞS v1 sin h1
w1 ¼ arctan ¼ arctan
ðv11 ÞP v1 cos h1 v
At Q2, the value of w2 results from
w2 ¼ w1 þ d
where d is the flyby turn-angle. As shown in the preceding figures, v2∞ rotates
counterclockwise with respect to v1∞ (i.e. the value of the angle d is positive), in
case of a leading-side flyby. By contrast, v2∞ rotates clockwise with respect to v1∞
(i.e. the value of the angle d is negative), in case of a trailing-side flyby.
Since the two vectors v1∞ and v2∞ have the same magnitude (v∞1 = v∞2 = v∞),
then the components of v2∞ along the directions uP and uS can be expressed as
follows
v21 ¼ ðv1 cos w2 Þ uP þ ðv1 sin w2 Þ uS
The velocity vector of the spacecraft at Q2 in its motion about the Sun is
v2 ¼ v21 þ v
The components of v2 along the directions uP and uS are
v2 ¼ ðv1 cos w2 þ vÞ uP þ ðv1 sin w2 Þ uS
Therefore, the radial component and the transverse component of v2 are
v2r ¼ v1 sin w2

v2/ ¼ v1 cos w2 þ v
The eccentricity, the true anomaly, and the moment of momentum per unit mass
of the spacecraft in its motion about the Sun at the point Q2 of exit from the sphere
of influence of the target planet are
h2 ¼ rP v2/ ¼ rP ðv1 cos w2 þ vÞ

2
h2 1
rP ¼
lS 1 þ e2 cos /2
l
v2r ¼ v1 sin w2 ¼ S e2 sin /2
h2
where rP (which is not to be confused with the radius of periapsis in the hyperbola)
is the distance of the spacecraft from the centre of mass of the Sun at the point Q2,
that is, in practice, the radius of the circular orbit of the target planet about the Sun.
As an example of application, we consider the flight past Jupiter of a spacecraft
coming from the Earth. The spacecraft leaves the orbit of the Earth at the perihelion
of its transfer orbit, which is a segment of ellipse having the Sun at a focus, and
encounters Jupiter at a point of this ellipse whose true anomaly is 160°. Let the
radius of periapsis of the spacecraft in its hyperbolic trajectory around Jupiter be rP
= 71492 + 1000 = 72492 km, where 71492 km is the radius of Jupiter. It is required
to compute the elements of the final elliptic orbit of the spacecraft about the Sun
after flyby. We consider this problem in two cases, depending on whether the
approach to Jupiter occurs on the dark side or on the illuminated side of the planet.
The data available at NASA [3, 7] are indicated below:
rE ¼ 1:496 1011 m
rJ ¼ 7:786 1011 m
lS ¼ 1:327 1020 m3 =s2
lJ ¼ 317:8 3:986 1014 m3 =s2
At the point of departure from the circular orbit of the Earth, which point is also
the perihelion of the elliptic segment from the Earth to Jupiter, there results
p1
rE ¼
1 þ e1
where p1 and e1 are, respectively, the semi-latus rectum and the eccentricity of the
elliptic segment of trajectory from the Earth to Jupiter.
The preceding equation, solved for p1, yields
p1 ¼ rE ð1 þ e1 Þ
At the point of intercept (/ = 160°) between this elliptic segment and the circular
orbit of Jupiter, the radius vector is
p1 r E ð 1 þ e1 Þ
rJ ¼ ¼
1 þ e1 cos 160 1 þ e1 cos 160
The preceding equation, solved for e1, yields
rJ rE 7:786 1:496
e1 ¼
¼ ¼ 0:7138
rE rJ cos 160 1:496 7:786 cos 160
Since the radius of perihelion of this elliptic segment is equal to the radius, rE, of
the circular orbit of the Earth, then the major semi-axis of this elliptic segment is
rE 1:496 1011
a1 ¼ ¼ ¼ 5:226 1011 m
1 e1 1 0:7138
The preceding expression, substituted into the vis-viva integral, makes it pos-
sible to compute the magnitude (in m/s) of the vector v1, which is the velocity of the
spacecraft with respect to the Sun at the point of intercept, as follows
1 12
2 1 2 2 1
v1 ¼ lS ¼ 1:327 10
20

rJ a1 7:786 1011 5:226 1011
¼ 9326
At the same point, the flight path angle results from

r0 e1 sin /
tan c1 ¼ ¼
r/0 1 þ e1 cos /

e1 sin / 0:7138 sin 160
c1 ¼ arctan ¼ arctan ¼ 36 :55
1 þ e1 cos / 1 þ 0:7138 cos160
The two components (transverse and radial) of the vector v1 are, respectively,
v1/ ¼ v1 cos c1 ¼ 9325:6 cos 36 :55 ¼ 7491 m=s

v1r ¼ v1 sin c1 ¼ 9325:6 sin 36 :55 ¼ 5554 m=s
As has been shown above, the two components of the velocity vector v1 along
the directions uP and uS can be expressed as follows
v1P ¼ v1/ ¼ 7491 m=s

v1S ¼ v1r ¼ 5554 m=s
Hence
v1 ¼ 7491 uP 5554 uS ðm=sÞ
On the other hand, the velocity vector of Jupiter at the intercept point is
12 1
lS 1:327 1020 2
v¼ uP ¼ uP ¼ 13055 uP ðm/s)
rJ 7:786 1011
Therefore, the hyperbolic excess velocity vector of the spacecraft, at the point of
entrance into the sphere of influence of Jupiter, is
v11 ¼ v1 v ¼ 7491 uP 5554 uS 13055 uP ¼ 5564 uP 5554 uS ðm=sÞ
The magnitude of this vector is

h i12
v1 ¼ ð5564Þ2 þ ð5554Þ2 ¼ 7861 m/s
The chosen radius of periapsis for the flyby hyperbola is
rP ¼ 72492 km ¼ 7:249 107 m
As has been shown in Sect. 5.6, the eccentricity and the semi-latus rectum of the
flyby hyperbola are, respectively,
rP v21 7:249 107 78612

e¼ þ1 ¼ þ 1 ¼ 1:035
lJ 317:8 3:986 1014
p ¼ rP ðe þ 1Þ ¼ 7:249 107 ð1:035 þ 1Þ ¼ 1:475 108 m
The angle between the transverse axis and the asymptote of the flyby hyperbola
results from

1 1
a ¼ arccos ¼ arccos ¼ 15 :02
e 1:035
The flyby turn-angle results from
d ¼ 180 2a ¼ 150 :0
For any elliptic orbit, the conjugate semi-axis results from
p 1:475 108
db¼ 1 ¼ 1 ¼ 5:499 108
ð e 2 1Þ 2
ð1:0352 1Þ2
The angle between v1∞ and v results from

v1S 5554
w1 ¼ arctan ¼ arctan ¼ 36 :55 ¼ 323 :4
v1P 7491
The spacecraft can approach the target planet in two different manners, as shown
in the following figure. When approaching Jupiter on the dark side, the spacecraft
turns counterclockwise along a branch of hyperbola, due to the gravitational force
exerted by the planet. The contrary is true when the approach occurs on the illu-
minated side of Jupiter.
Let us consider the case of an approach on the dark side of Jupiter. In this case,
the value of the flyby turn-angle is positive (d = 150°.0), because the spacecraft
turns counterclockwise, as shown in the following figure.
Since the values of the angles w1 and d have been determined above, we can
compute the value of the angle w2 between the hyperbolic excess velocity vector
v2∞ of the spacecraft at the point of exit from the sphere of influence of Jupiter and
the velocity vector v of Jupiter with respect to the Sun at the same point, as follows
w2 ¼ w1 þ d ¼ 36 :55 þ 150 :0 ¼ 113 :4
Since the value of w2 has been computed, we can determine the vector v2∞ at the
point of exit as follows
where v∞ = 7861 m/s. Hence
v21 ¼ ð7861 cos 113 :4Þ uP þ ð7861 sin 113 :4Þ uS ¼ 3123 uP þ 7214 uS
The velocity vector of the spacecraft with respect to the Sun at the same point
results from
v2 ¼ v21 þ v ¼ 3123 uP þ 7214 uS þ 13055 uP ¼ 9932 uP þ 7214 uS
where v = 13055 uP is the velocity vector of Jupiter at the same point.

The components of v2 along, respectively, the transverse direction and the radial
direction are
v2/ ¼ 9932 m=s

v2r ¼ 7214 m=s
The magnitude of the vector v2 is

12 h i12
v2 ¼ v22/ þ v22r ¼ 99322 þ ð7214Þ2 ¼ 12275 m/s
This shows an increase v2 − v1 = 2949 m/s after flyby with respect to v1 =

9326 m/s before flyby. This is because of the vector sum v2 = v2∞ + v. Therefore,
the spacecraft accelerates after flyby.
The moment of momentum per unit mass of the spacecraft after flyby results
from
h2 ¼ rJ v2/ ¼ 7:786 1011 9932 ¼ 7:733 1015 m2 =s
Hence, the semi-latus rectum of the elliptic orbit after flyby is

2
h22 7:733 1015
p2 ¼ ¼ ¼ 4:506 1011 m
lS 1:327 1020
The general equation of the conic sections

p
r¼
1 þ e cos /
yields
p2 4:506 1011
e cos /2 ¼ 1¼ 1 ¼ 0:4213
rJ 7:786 1011
On the other hand, the equation

lS
v2r ¼ e2 sin /2
h2
solved for e2 sin /2 yields
v2r h2 7214 7:733 1015

e2 sin /2 ¼ ¼ ¼ 0:4204
lS 1:327 1020
Hence

e2 sin /2 0:4204
/2 ¼ arctan ¼ arctan ¼ 224 :9
e2 cos /2 0:4213
This is because sin /2 and cos /2 have, both of them, negative values.
Since /2 = 224°.9 and e2 cos /2 = −0.4213, then
0:4213
e2 ¼ ¼ 0:5951
cos 224 :9
This is the eccentricity of the elliptic orbit of the spacecraft around the Sun after
flyby. The true anomaly /2 = 224°.9 specifies the position of the spacecraft in its
elliptic orbit around the Sun after flyby. Since the point of intercept between the
transfer ellipse and the orbit of Jupiter has been chosen at an angular distance of
160° from the point of departure, then the perihelion of the elliptic orbit after flyby
is placed at an angular distance of 224°.9 − 160°.0 = 64°.9 counterclockwise from
the point of intercept, and the aphelion of the same orbit is placed at an angular
distance of 64°.9 + 180°.0 = 244°.9 counterclockwise from the point of intercept.
The radius of perihelion of the elliptic orbit after flyby results from setting / = 0°
in the equation r = p2/(1 + e2 cos /). This yields
4:506 1011
rPH ¼ ¼ 2:825 1011 m
ð1 þ 0:5951Þ
This value is less than that (rJ = 7.786 1011 m) of the radius of the circular
orbit of Jupiter. The radius of aphelion of the elliptic orbit after flyby results from
setting / = 180° in the equation r = p2/(1 + e2 cos /). This yields
4:506 1011
rAH ¼ ¼ 1:113 1012 m
ð1 0:5951Þ
This value is greater than that (rJ = 7.786 1011 m) relating to the radius of the
circular orbit of Jupiter.
The total energy per unit mass of the spacecraft before the encounter with Jupiter
is
lS 1:327 1020
E1 ¼ ¼ ¼ 1:270 108 m2 =s2
2a1 2 5:226 1011
The total energy per unit mass of the spacecraft after the encounter with Jupiter
is
lS 1:327 1020
E2D ¼ ¼ ¼ 9:509 107 m2 =s2
rPH þ rAH 2:825 1011 þ 1:113 1012
Since the value of E2D is greater than the value (−1.270 108 m2/s2) of E1, then
the spacecraft has gained energy after flyby by subtracting it from the energy
possessed by Jupiter.
Let us consider now the case of an approach on the illuminated side of Jupiter. In
this case, the value of the flyby turn-angle is negative (d = −150°.0), because the
spacecraft turns clockwise, as has been shown in the preceding figure.
The value of the angle w2 between the hyperbolic excess velocity vector v2∞ of
the spacecraft at the point of exit from the sphere of influence of Jupiter and the
velocity vector v of Jupiter with respect to the Sun at the same point results from
w2 ¼ w1 þ d ¼ 323 :4 150 :0 ¼ 173 :5
The vector v2∞ at the point of exit has the following components along uP and
uS
where v∞ = 7861 m/s. Hence
v21 ¼ ð7861 cos 173 :5Þ uP þ ð7861 sin 173 :5Þ uS ¼ 7811 uP þ 892:0 uS
The velocity vector of the spacecraft with respect to the Sun at the same point
results from
v2 ¼ v21 þ v ¼ 7811 uP þ 892 uS þ 13055 uP ¼ 5244 uP þ 892:0 uS
where v = 13055 uP is the velocity vector of Jupiter at the same point.

The components of v2 along, respectively, the transverse direction and the radial
direction are
v2/ ¼ v2P ¼ 5244 m=s

v2r ¼ v2S ¼ 892:0 m=s
The magnitude of the vector v2 is

12 h i12
v2 ¼ v22/ þ v22r ¼ 52442 þ ð892:0Þ2 ¼ 5320 m/s
This shows a decrease v1 − v2 = 4006 m/s after flyby with respect to v1 =

9326 m/s before flyby. Therefore, the spacecraft decelerates after flyby.
The moment of momentum per unit mass of the spacecraft after flyby is
h2 ¼ rJ v2/ ¼ 7:786 1011 5244 ¼ 4:083 1015 m2 =s
Hence, the semi-latus rectum of the elliptic orbit after flyby is

2
h22 4:083 1015
p2 ¼ ¼ ¼ 1:257 1011 m
lS 1:327 1020
The quantities e2 cos /2 and e2 sin /2 are, respectively,
p2 1:257 1011
e cos /2 ¼ 1¼ 1 ¼ 0:8386
rJ 7:786 1011
v2r h2 892:0 4:083 1015
e2 sin /2 ¼ ¼ ¼ 0:02745
lS 1:327 1020
Hence

e2 sin /2 0:02745
/2 ¼ arctan ¼ arctan ¼ 181 :9
e2 cos /2 0:8386
This is because sin /2 and cos /2 have, both of them, negative values.
Since /2 = 181°.9 and e2 cos /2 = −0.8386, then
0:8386
e2 ¼ ¼ 0:8391
cos 181 :9
This is the eccentricity of the elliptic orbit of the spacecraft around the Sun after
flyby. The true anomaly /2 = 181°.9 specifies the position of the spacecraft in its
elliptic orbit around the Sun after flyby. Since the point of intercept between the
transfer ellipse and the orbit of Jupiter has been chosen at an angular distance of
160° from the point of departure, then the perihelion of the elliptic orbit after flyby
is placed at an angular distance of 181°.9 − 160°.0 = 21°.9 counterclockwise from
the point of intercept, and the aphelion of the same orbit is placed at an angular
distance of 21°.9 + 180°.0 = 201°.9 counterclockwise from the point of intercept.
The radius of perihelion of the elliptic orbit after flyby results from setting / = 0
in the equation r = p2/(1 + e2 cos /). This yields
1:257 1011
rPH ¼ ¼ 6:832 1010 m
ð1 þ 0:8391Þ
This value is less than that (rJ = 7.786 1011 m) of the radius of the circular
orbit of Jupiter. The radius of aphelion of the elliptic orbit after flyby results from
setting / = 180° in the equation r = p2/(1 + e2 cos /). This yields
1:257 1011
rAH ¼ ¼ 7:808 1011 m
ð1 0:8391Þ
This value is greater than that (rJ = 7.786 1011 m) relating to the radius of the
circular orbit of Jupiter.
The total energy per unit mass of the spacecraft after the encounter with Jupiter
is
lS 1:327 1020
E2I ¼ ¼ ¼ 1:563 108 m2 =s2
rPH þ rAH 6:832 1010 þ 7:808 1011
Since the value of E2I is less than the value (−1.270 108 m2/s2) of E1, then the
spacecraft has lost energy after flyby in favour of Jupiter.
The increase (or decrease) in the magnitude of the velocity vector of a spacecraft
with respect to the Sun after the encounter with a target planet can be used to accelerate
(or decelerate) the spacecraft in its path, as will be shown in the next paragraph.
5.9 The Gravity Assist
The gravity assist (also known as the slingshot effect) is a technique used to
increase or decrease the velocity of a spacecraft with respect to the Sun after an
encounter between that spacecraft and an assisting planet.
By so doing, momentum is added or subtracted to increase or decrease the

mechanical energy of a spacecraft in its orbit about the Sun. This fact does not
violate the fundamental law of conservation of energy, because the amount of
energy gained by the spacecraft is subtracted from the energy possessed by the
assisting planet in their respective orbits about the Sun. The gravity assist technique
has been used in several space missions. For example, according to NASA-JPL [8],
the Voyager 2 spacecraft was launched in August 1977 and flew by Jupiter, for the
purpose of taking gravity assist for a trajectory boost to Saturn. Voyager 1 was
launched the following month and did the same, then obtained a gravity assist from
Saturn and a further one from Uranus, so as to reach Neptune and go beyond. The
Galileo spacecraft used gravity assist to get a boost from Venus and two from the
Earth, in order to reach Jupiter. The Cassini–Huygens spacecraft, whose trajectory
is shown in the preceding figure, due to the courtesy of NASA-JPL [9], took two
gravity assists from Venus, one from the Earth and another from Jupiter to gain the
momentum required to reach Saturn.
A gravity assist involves three bodies. They are a spacecraft, a celestial body
(e.g. Venus, the Earth, Jupiter, …) whose energy is extracted in favour of the
spacecraft, and another body (usually the Sun), about which the path of the
spacecraft is to be modified.
In a model based on two bodies (e.g. Voyager 2 interacting only with Jupiter),
there is an exchange of kinetic and potential energy between the spacecraft and the
assisting planet. The spacecraft increases its kinetic energy and decreases its potential
energy when approaching to the planet, and does the contrary when leaving the planet
after flyby. There is no gain or loss of velocity of one body with respect to the other,
and the total energy is constant throughout the whole process. As has been shown in
Sect. 5.8, the gravitational attraction of the planet produces the only effect of rotating
the hyperbolic excess velocity vector, v1∞, of the spacecraft at the point of entrance
into the sphere of influence of the planet. As a result of this rotation, the hyperbolic
excess velocity vector, v2∞, of the spacecraft at the point of exit from the sphere of
influence of the planet has a different direction with respect to v1∞, but the same
magnitude, as shown in the following figure, due to the courtesy of NASA-JPL [8].
5.9 The Gravity Assist 813
In the real world, the spacecraft interacts not only with the assisting planet, but
also with the Sun, because this planet moves with respect to the Sun, as shown in
the following figure, also due to the courtesy of NASA-JPL [8].
Such being the case, an exchange of energy occurs between the spacecraft and
the planet which leaves the total energy of each of them unaltered with respect to
the total energy of the other, but changes the energy of each of them with respect to
the Sun. When the spacecraft gains energy from the assisting planet, then the planet
loses energy, and vice versa, depending on the type of trajectory followed by the
spacecraft. However, the total energy possessed by the three bodies (the spacecraft,
the assisting planet, and the Sun) does not change, in accordance with the law of
conservation of energy. Since the momentum of a body depends on the mass and
velocity of that body, and since the masses of the bodies involved in a gravity assist
remain unchanged, then only the velocities of those bodies change. Due to the law
of conservation of momentum, the assisting planet, whose mass is much greater
than that of the spacecraft, imparts a significant velocity to the spacecraft at the cost
of a negligible change in its own velocity.
Let mS and mP be the masses of, respectively, the spacecraft and the assisting
planet. Let vS and vP be the respective velocity vectors. In the following expres-
sions, the subscript 1 and 2 indicate these vectors, respectively, before and after
flyby. The law of conservation of momentum requires that
mS vS1 þ mP vP1 ¼ mS vS2 þ mP vP2
Hence
mS
vP2 vP1 ¼ ðvS1 vS2 Þ
mP
The mass ratio, mS/mP, is usually very small, because mS is about 103 kg and mP
ranges from 0.330 1024 (Mercury) to 1898 1024 (Jupiter), and consequently,
mS/mP ranges from about 10−20 to about 10−24. Consequently, to a very small
decrease in the velocity of the assisting planet, there corresponds a very high
increase in the velocity of the spacecraft. As has been shown in Sect. 5.8, a gravity
assist increases the velocity of a spacecraft when the vector vS2∞ − vS1∞ has a
positive component along the direction of the velocity vector vP of the assisting
planet with respect to the Sun. By contrast, a gravity assist decreases the velocity of
a spacecraft when the vector vS2∞ − vS1∞ has a negative component along the
direction of the velocity vector vP of the assisting planet with respect to the Sun.
Johnson [10] has shown the range of velocities vS2, which a spacecraft can reach
after gravity assist with a given velocity vS1 of the same spacecraft before gravity
assist and a given orbital velocity vP of the assisting planet, as follows:
• the maximum increase of velocity for a spacecraft in a gravity assist is achieved
when vS2 aligns exactly with vP;
• any other increase of velocity than the maximum can be achieved from either of
two gravity assists, when the corresponding vectors vS2 are related by reflection
in vP;
• too large a deflection brakes a spacecraft, that is, vS2 < vS1;
• the maximum possible gain is 2vP, for a head-on collision with rotation through
180°.
The braking effect obtainable by means of a gravity assist can be used to save fuel
for the purpose of inserting a spacecraft into a capture orbit around a planet. Such has
been the case with the Cassini spacecraft for an orbit insertion about Saturn. In addition,
the Galileo spacecraft has braked at Jupiter with the help of the satellite Io [10].
A gravity assist provided by Jupiter has also been used in the joint ESA-NASA
Ulysses mission to deflect the path of a spacecraft out of the ecliptic plane, as
shown in the following figure, due to the courtesy of NASA [11].
5.9 The Gravity Assist 815
A procedure to perform gravity assist calculations and an account of the Voyager

1 Jupiter flyby have been provided by Cesarone [12]. A tutorial on gravity assist in
celestial mechanics is due to Van Allen [13].
5.10 Orbital Elements of the Planets
For the purpose of computing an interplanetary trajectory, it is necessary to

determine the orbital elements (or, which is the same, the position and velocity
vectors) of the planets involved at the time of interest. There are several methods to
do this. Two of them are shown below.
Meeus [14] indicates a method which expresses the orbital elements of the
planets from Mercury to Neptune as polynomials having the following form
a0 þ a1 T þ a2 T 2 þ a3 T 3
where T is the time measured in Julian centuries of 36525 ephemeris days from the
epoch J2000.0 = 1.5 January 2000 = JDE 2451545.0.
In other words,
JDE 2451545:0
T¼
36525
The concept of Julian ephemeris day (JDE) has been given to the reader in
Sect. 3.4.
The quantity T resulting from this calculation is negative before the beginning of
the year 2000, and positive afterwards. The orbital elements of the planets are
L = mean longitude of the planet;
a = major semi-axis of the elliptic orbit of the planet;
e = eccentricity of the orbit;
i = inclination of the orbit with respect to the plane of the ecliptic;
X = longitude of the ascending node;
P = X + x = longitude of the perihelion.
Some authors use the letter - to indicate the longitude of the perihelion. We use
the letter P in order to avoid confusion with the argument of the perihelion (x).
The quantities L and P are measured in two different planes, namely from the
vernal equinox along the ecliptic to the ascending node of the orbit and then from
the ascending node along the orbit. The mean anomaly of a planet is given by M =
L − P. Meeus [14] gives the values of the coefficients a0, a1, a2, and a3 in a table
(Table 30.A, pp. 200–202) which refers to the ecliptic and the mean equinox of the
date. The values for the major semi-axis are expressed in astronomical units (AU).
The values for the angles L, i, X, and P are expressed in degrees and decimals.
These values can also be found in Ref. [15].
In some cases, it is desirable to refer the elements L, i, X, and P to a standard

equinox. In such cases, it is possible to determine the values of the four coefficients
named above by means of another table (Table 30.B, pp. 203–204) which refers to
the standard equinox of J2000.0. The values of the coefficients for the orbital
elements a and e, which are not modified by a change of reference system, are to be
computed by means of Table 30.A. For the Earth, in order to avoid a discontinuity
in the variation of the inclination and a jump of 180° in the longitude of the
ascending node at the epoch J2000.0, the inclination with respect to the ecliptic of
J2000.0 is considered negative before the year 2000.
As an example, let us compute the orbital elements of, respectively, Mars and the
Earth on 8 March 2012, noon UT1. As has been shown in Sect. 3.14, the corre-
sponding value of T is 0.121834381476.
By using the values given in Table 30.A of Ref. [14], we find for Mars
L ¼ 355:433275 þ 19141:6964746 T þ 0:00031097 T 2 þ 0:000000015 T 3

¼ 167 :55003
a ¼ 1:523679342 AU
e ¼ 0:09340062 þ 0:000090483 T 0:0000000806 T 2 0:00000000035 T 3
¼ 0:09341164274331
i ¼ 1:849726 0:0006010 T þ 0:00001276 T 2 0:000000006 T 3
¼ 1 :84965296693
X ¼ 49:558093 þ 0:7720923 T þ 0:00001605 T 2 þ 0:000002325 T 3
¼ 49 :65216063
P ¼ 336:060234 þ 1:8410331 T þ 0:00013515 T 2 þ 0:000000318 T 3
¼ 336 :2845371357
The value of L computed above has been reduced to the range 0°–360°.
Likewise, we find for the Earth
L ¼ 100:466449 þ 36000:7698231 T þ 0:00030368 T 2 þ 0:000000021 T 3

¼ 166 :59797756498
a ¼ 1:000001018 AU
e ¼ 0:01670862 0:000042037 T 0:0000001236 T 2 þ 0:00000000004 T 3
¼ 0:0167034966135
i ¼ 0 :0
P ¼ 102:937348 þ 1:7195269 T þ 0:00045962 T 2 þ 0:000000499 T 3
¼ 103 :146852319618
5.10 Orbital Elements of the Planets 817
The value of L computed above has been reduced to the range 0°–360°. When
the orbital elements are computed for the Earth, the value of X is not defined,
because i = 0°.0. This value may arbitrarily be set to X = 0°.0.
Another method for computing the orbital elements of the planets is based on the
following table, due to the courtesy of NASA-JPL [16].
This table is to be used as specified by Standish [17]. In brief, the table gives the
orbital elements of the planets and their rates of change per century in the units
specified above. The elements a, e, i, and L are the same as those defined above.
The following column indicates the longitude of the perihelion (P) and its rate of
change per century. The rightmost column indicates the longitude of the ascending
node (X) and its rate of change per century.
In order to use the second method, the Julian Day Number (J0) and the Julian
Date (JD) are computed as specified in Sect. 3.14. For convenience of the reader,
with reference to the preceding example (8 March 2012, noon UT1), we set y =
2012, m = 3, d = 8, and h = 12. Then, J0 and JD result from
J0 ¼ 367y INTf1:75½y þ INTðm=12 þ 0:75Þ

g þ INTð275m=9Þ þ d þ 1721013:5
¼ 367 2012 INTf1:75 ½2012 þ INTð3=12 þ 0:75Þ
g þ INTð275 3=9Þ þ 8
þ 1721013:5 ¼ 2455994:5
JD ¼ J0 þ h=24 ¼ 2455994:5 þ 12=24 ¼ 2455995:0
The time measured in Julian centuries results from
JD 2451545:0 2455995:0 2451545:0

T0 ¼ ¼ ¼ 0:121834360027
36525 36525
For each planet, the upper row of the preceding table specifies the coefficient a0
and the lower row specifies the coefficient a1 of a first-degree polynomial
a0 þ a1 T0
which is a function of the time T0. This holds for each column of the table, which
relates to an individual orbital element.
For example, the orbital elements of Mars at T0 = 0.121834360027 are computed
as follows:
Mean longitude:
L ¼ 4:55343205 þ 19140:30268499 0:121834360027 ¼ 167 :3930962988
Again, the value of L computed above has been reduced to the range 0°–360°.
This value is to be compared with L = 167°.55003 computed previously by means
of the method described by Meeus [14].
Major semi-axis:
a ¼ 1:52371034 þ 0:00001847 0:121834360027 ¼ 1:52371259028 AU
Eccentricity:
e ¼ 0:09339410 þ 0:00007882 0:121834360027 ¼ 0:093403702984257
Inclination:
i ¼ 1:84969142 0:00813131 0:121834360027 ¼ 1 :8487007450964785
Longitude of the ascending node:
X ¼ 49:55953891 0:29257343 0:121834360027 ¼ 49 :523893413395
Longitude of the perihelion:
P ¼ 23:94362959 þ 0:44441088 0:121834360027 ¼ 336 :1105149251538
When the orbital elements for the planets of interest have been computed by
means of one of the two methods shown above, the positions of these planets can
also be expressed in heliocentric ecliptic rectangular co-ordinates X, Y, and Z, which
are shown in the following figure.
The origin of these co-ordinates is the centre of mass of the Sun. The funda-
mental plane, XY, is the plane of the ecliptic. The X-axis points towards the vernal
equinox of some reference date (e.g. J2000.0). The Y-axis is 90° to the east of the X-
axis in the fundamental plane. Finally, the Z-axis is perpendicular to the XY plane
and points towards the north ecliptic pole, such that uZ = uX uY.
In order to convert the orbital elements of a given planet into heliocentric ecliptic
co-ordinates, it is necessary to convert the major semi-axis, which is expressed in
AU, into the desired unit of length. For example, in case of metric units, 1 AU is
exactly equal to 149597870700 m.
The major semi-axis, a, and the eccentricity, e, relating to the orbit of the planet
of interest at the given time are used to compute the moment of momentum per unit
mass, as follows
1 1
h ¼ ðlS pÞ2 ¼ lS a 1 e2 2
where lS = 1.327 1020 m3/s2 is the gravitational parameter of the Sun.

Instead of h, it is also possible to compute the semi-latus rectum
p ¼ að1 e2 Þ
The argument of perihelion, x, and the mean anomaly at epoch, M, result from
x¼PX
M ¼LP
By using the values of e and M, with the latter expressed in radians, Kepler’s
equation M = Æ − e sin Æ is solved iteratively for the eccentric anomaly.
As has been shown in Sect. 1.3, the true anomaly, /, results from
1
/ 1þe 2 Æ
tan ¼ tan
2 1e 2
By using the classical elements computed above, it is possible to compute the

position, r, and velocity, v, vectors for the planet of interest at the given time, as has
been shown in Sect. 1.9.
For example, let us compute the state vector of Mars on 8 March 2012, at noon
UT1.
We use the following set of orbital elements
L ¼ 167 :393096 i ¼ 1 :84870075

a ¼ 1:52371259 AU X ¼ 49 :5238934
e ¼ 0:0934037030 P ¼ 336 :110515
which have been computed above by means of the second method.

The major semi-axis of the orbit of Mars, expressed in metres, is
a ¼ 1:52371259 1:495978707 1011 ¼ 2:27944563 1011 m
The semi-latus rectum is
p ¼ að1 e2 Þ ¼ 2:27944563 1011 ð1 0:09340370302 Þ

¼ 2:25955917 1011 m
The argument of the perihelion is
x ¼ P X ¼ 336 :110515 49 :5238934 ¼ 286 :586622
The mean anomaly at epoch is
M ¼ L P ¼ 167 :393096 336 :110515 ¼ 168 :717419
The same value, reduced to the range 0°–360° and then converted into radians, is
M = 3.33851084. A starting value of eccentric anomaly, used to solve iteratively
Kepler’s equation, is provided by the following expression suggested by Battin [18,
p. 194]:
e sin M
Æ0 ¼ M þ
1 sinðM þ eÞ þ sin M
0:0934037030 sin 3:33851084

Æ 0 ¼ 3:33851084 þ
1 sinð3:33851084 þ 0:0934037030Þ þ sin 3:33851084
¼ 3:32175489
The Newton–Raphson method (see Sect 1.5) is used as follows
Mk ¼ Æ k e sin Æ k
M Mk
Æk þ 1 ¼ Æk þ
1 e cos Æ k
By so doing, we have
Æ 0 ¼ 3:32175489
M0 ¼ 3:32175489 0:0934037030 sin 3:32175489 ¼ 3:3384918
3:33851084 3:3384918
Æ 1 ¼ 3:32175489 þ ¼ 3:32177233
1 0:0934037030 cos 3:32175489
M1 ¼ 3:32177223 0:0934037030 sin 3:32177223 ¼ 3:33851086
3:33851084 3:33851086
Æ 2 ¼3:32177233 þ ¼ 3:32177231
1 0:0934037030 cos 3:32177233
M2 ¼ 3:32177231 0:0934037030 sin 3:32177231 ¼ 3:33851084
Since M2 is equal to M within the chosen tolerance, then Æ2 = 3.32177231 is

accepted as the solution of Kepler’s equation. The true anomaly at epoch, /0,
results from
1 1
/0 1þe 2 Æ 1 þ 0:093403703 2 3:32177231
tan ¼ tan ¼ tan
2 1e 2 1 0:093403703 2
¼ 12:1571108
Hence
/0 ¼ 2 arctan ð12:1571108Þ ¼ 2:97744942 rad ¼ 189 :404715
The orbital elements of Mars computed above at the given epoch (8 March 2012,
noon UT1), that is
a ¼ 2:27944563 1011 m e ¼ 0:0934037030 i ¼ 1 :84870075

X ¼ 49 :5238934 x ¼ 286 :586622 /0 ¼ 189 :404715
are used to compute the three components X, Y, and Z of the position vector r and
the three components VX, VY, and VZ of the velocity vector r′ of Mars in the
heliocentric ecliptic reference system (ICRF, J2000.0), at epoch t0. By using the
method shown in Sect. 1.9, we obtain
X ¼ 2:40890556 1011 m Y ¼ 6:21775171 1010 m Z ¼ 7:21720117 109 m

VX ¼ 5:14897261 103 m=s VY ¼ 2:13916561 104 m=s VZ ¼ 3:21779514 102 m=s
For comparison, the results obtained by means of the web interface

NASA-JPL-HORIZONS [19] at the same epoch, expressed in km and km/s, are
also given below.
5.11 General Interplanetary Trajectories
In the general case, the interplanetary trajectory of a spacecraft, which goes from
the sphere of influence of a planet of departure to the sphere of influence of a planet
of arrival, does not lie on the ecliptic plane (i.e. on the XY plane of the heliocentric
ecliptic reference system XYZ defined in Sect. 5.10), nor is this trajectory a
Hohmann semi-ellipse. It is a segment of a conic section having the centre of mass
of the Sun at one of its foci.
In order to determine an interplanetary trajectory, it is necessary first of all to
determine the state vector of the planet of departure at the time, td, of departure and
the state vector of the planet of arrival at the time, ta = td + Dt, of arrival.
The methods which can be used to determine the state vector of a given planet at
a given time have been shown in Sect. 5.10.
This done, it is necessary to determine the trajectory of the spacecraft from the
planet of departure to the planet of arrival in the heliocentric ecliptic reference
system XYZ. For this purpose, the position vector of the spacecraft at the time, td, of
departure is the same as that of the planet of departure. Likewise, the position vector
of the spacecraft at the time, td + Dt, of arrival is the same as that of the planet of
arrival. When these two position vectors and the time, Dt, of transfer are specified
by the designer, it is possible to determine the two velocity vectors of the spacecraft
at the times of, respectively, departure and arrival by solving the problem of
Lambert (i.e. the determination of an orbit from two position vectors, the time of
flight, and the direction of flight), as shown in Sects. 1.5 and 1.8.
Either the state vector of the spacecraft at the time td or the state vector of the
spacecraft at the time td + Dt can be used to determine the six orbital elements
which identify the transfer trajectory of the spacecraft between the spheres of
influence of the two planets of, respectively, departure and arrival.
According to the nomenclature of Sect. 5.5, let vd be the velocity vector of the
spacecraft in the heliocentric ecliptic reference system XYZ at the time, td, of leaving
5.11 General Interplanetary Trajectories 823
the sphere of influence of the planet of departure. The hyperbolic excess velocity
vector of the spacecraft with respect to the planet of departure at t = td results from
vd1 ¼ vd vPd
where vPd is the orbital velocity of the planet of departure with respect to the Sun at
t = td. The magnitude of this hyperbolic excess velocity vector is
vd1 ¼ jvd vPd j
Likewise, let va be the velocity vector of the spacecraft in the heliocentric ecliptic
reference system XYZ at the time, ta = td + Dt0, of reaching the sphere of influence of
the planet of arrival. The hyperbolic excess velocity vector of the spacecraft with
respect to the planet of arrival at t = ta results from
va1 ¼ va vPa
where vPa is the orbital velocity of the planet of arrival with respect to the Sun at
time t = ta. The magnitude of this hyperbolic excess velocity vector is
va1 ¼ jva vPa j
The procedure to be followed to compute the transfer trajectory between the

spheres of influences of two planets is described by means of the following
example, which has been proposed by Curtis [5].
A spacecraft leaves the sphere of influence of the Earth, on 7 November 1996, at
00:00:00 UTC, directed in a prograde trajectory to Mars. This spacecraft reaches the
sphere of influence of Mars on 12 September 1997, at 00:00:00 UTC. The time of
transfer resulting from the departure and arrival days chosen is Dt =
2450703.5 – 2450394.5 = 309 days = 2.66976 107 s.
We want to compute the interplanetary trajectory of the spacecraft between these
two spheres of influence in the heliocentric ecliptic reference system (ICRF,
J2000.0). To this end, it is necessary to determine the position and velocity vectors
for, respectively, the Earth at the time of departure and Mars at the time of arrival.
For this purpose, it is possible to use the methods described in Sect. 5.10. However,
it is easier to determine these vectors by means of the web interface
NASA-JPL-HORIZONS [19].
The results found in Ref. [19] are given below.
Earth at the time of departure:
Mars at the time of arrival:
In the two sets of data given above, the three rows indicate, respectively, the
Julian date, the three components X, Y, and Z of the position vector of the planet in
km, and the three components VX, VY, and VZ of the velocity vector of the planet in
km/s. The position vectors of the spacecraft in the heliocentric ecliptic reference
system XYZ at the times of, respectively, departure (subscript d) from the Earth and
arrival (subscript a) at Mars are the same as the position vectors of, respectively, the
Earth and Mars. By using nine decimal places, these vectors are
rd ¼ 1:04998587 108 uX þ 1:04650719 108 uY þ 1:12120077 103 uZ

ra ¼ 2:08489513 107 uX 2:18416657 108 uY 4:06276688 106 uZ
where uX, uY, and uZ are unit vectors relating to, respectively, X, Y, and Z.
The position vectors rd and ra specified above and the chosen time of transfer Dt
are used to compute the two velocity vectors vd and va of the spacecraft in the
heliocentric ecliptic reference system XYZ at the times of, respectively, departure
and arrival by solving the problem of Lambert, as will be shown below.
The gravitational parameter of the Sun is lS = 1.32712440018 1011 km3/s2;
1
hence, the square root of this value is lS =2 = 3.64297186 105 km3/2/s. The
magnitudes of the position vectors rd and ra are, respectively,
h 2 2 2 i12
rd ¼ 1:04998587 108 þ 1:04650719 108 þ 1:12120077 103
¼ 1:48244650 108 km
h 2 2 2 i12
ra ¼ 2:08489513 107 þ 2:18416657 108 þ 4:06276688 106
¼ 2:19447080 108 km
Let D/ be the angle through which the spacecraft travels along its interplanetary
trajectory from the Earth to Mars. The cosine of this angle results from
rd ra
cos D/ ¼
rd ra
By substituting the components and the magnitudes of the vectors rd and ra into
the preceding equation, we obtain

rd ra ¼ 1:04998587 108 2:08489513 107 þ 1:04650719 108

2:18416657 108 þ 1:12120077 103

4:06276688 106 ¼ 2:50465731 1016 km2
rd ra ¼ 1:48244650 108 2:19447080 108 ¼ 3:25318556 1016 km2
rd ra 2:50465731 1016
cos D/ ¼ ¼ ¼ 0:76990921
rd ra 3:25318556 1016
Since cos D/ < 0, then the transfer angle D/ may be either in the second
quadrant or in the third quadrant. In order to determine the correct quadrant for D/,
we compute the component of the vector product rd ra along the Z-axis of the
heliocentric ecliptic reference system XYZ. In the present case, the planetary
co-ordinates at departure and arrival are such that

ðrd ra ÞZ ¼ Xd Ya Yd Xa ¼ 1:04998587 108 2:18416657 108

1:04650719 108 2:08489513 107 ¼ 2:07515826 1016 km2
Since we want the spacecraft to go from the Earth to Mars along a prograde
trajectory (such that 0 < i < p/2) and (rd ra)Z is less than zero, then the transfer
angle, D/, of the spacecraft is in the third quadrant; that is, there results

rd ra
D/ ¼ 2p arccos ¼ 2p arccosð0:76990921Þ ¼ 3:83369011 rad
rd ra
¼ 219 :654263
Generally speaking, when
ðrd ra ÞZ 0
then

rd ra
D/ ¼ arccos for a prograde orbit
rd ra
rd ra
D/ ¼ 2p arccos for a retrograde orbit
rd ra
By contrast, when
ðrd ra ÞZ \0
then

rd ra
D/ ¼ 2p arccos for a prograde orbit
rd ra
rd ra
D/ ¼ arccos for a retrograde orbit
rd ra
The method used below to solve the problem of Lambert relating to the present
case is the same as that described in Sect. 1.7.
Since the transfer angle D/ is in the third quadrant, then the minus sign takes
effect in front of the following expression
1
A ¼ ½rd ra ð1 þ cos D/Þ
2
and therefore
1
A ¼ 1:48244650 108 2:19447080 108 ð1 0:76990921Þ 2
¼ 8:65175147 107 km
The lower endpoint (zL) and the upper endpoint (zU) of the interval, within which
the unknown value of the variable z is sought, are taken initially as follows
zL ¼ 4p2 39:4
zU ¼ 4p2 39:4
First, in order to determine the sign of the variable y at the point z = zU, we
compute the values of the two Stumpff functions C(z) and S(z) at this point, as
follows

1

1 cos 39:42
CU ¼ Cð39:4Þ ¼ ¼ 4:94665287 107
39:4

1 1
39:42 sin 39:42
SU ¼ Sð39:4Þ ¼ 3 ¼ 0:0254059555
39:42
The correspondent value of y, at z = zU, is
AðzU SU 1Þ
yU ¼ r d þ r a þ 1 ¼ 1:48244650 108 þ 2:19447080 108
2
CU
39:4 0:0254059555 1
8:65175147 107 1 ¼ 2:45338014 108 km
ð4:94665287 107 Þ2
Therefore, at z = zU, there results yU > 0. Then, zU = 39.4 is a possible value of z,

and we can go further in the computation, as follows
12
2:45338014 108 1
xU ¼ ¼ 2:22703327 107 km2
4:94665287 107
function
1
x3U SU þ A y2U
f ðDt; xU Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, xU = 2.22703327 107, SU =

1
0.0254059555, A = −8.65175147 107, yU = 2.45338014 108, and lS =2 =
3.64297186 10 into the preceding equation, we find
5
f ðDt; xU Þ ¼ 7:70299853 1014 s
The value of f(Dt, xU) is negative and very far from zero. Consequently, at z = zU,
we are largely away from the point x in which the function f(Dt, x) vanishes. The
sign in front of f(Dt, xU) and the sign in front of f(Dt, xL) are important, because, in
case of f(Dt, x) having opposite signs at the endpoints xL and xU, there is assurance
of the existence of a zero of f(Dt, x) between xL and xU.
Now, in order to determine the sign of y at the lower endpoint (zL = −39.4) of the
interval, we compute the values of the Stumpff functions C(z) and S(z) at z = zL, as
follows

1

1 cosh 39:42
CL ¼ C ð39:4Þ ¼ ¼ 6:72792733

39:4
1
1
sinh 39:42 39:42
SL ¼ Sð39:4Þ ¼ 3 ¼ 1:05050312
39:42
The correspondent value of y, at z = zL, is
AðzL SL 1Þ
yL ¼ rd þ ra þ 1 ¼ 1:48244650 108 þ 2:19447080 108
CL
2
39:4 1:05050312 1
8:65175147 107 1 ¼ 1:78161223 109 km
6:727927332
Thus, at z = zL, there results yL > 0. Then, zL = −39.4 is also a possible value of z,
and we can go further in the computation, as follows
12 12
yL 1:78161223 109 1
xL ¼ ¼ ¼ 16272:9371 km2
CL 6:72792733
At x = xL and within the interval zL < z < zU, we seek a zero of the function
1
x3L SL þ A y2L
f ðDt; xL Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, xL = 16272.9371, SL = 1.05050312, A =

1
−8.65175147 107, yL = 1.78161223 109, and lS =2 = 3.64297186 105 into
the preceding equation, we find
f ðDt; xL Þ ¼ 2:42956879 107 s
The value of the time function f(Dt, xL) is positive. Consequently, there is
assurance of the existence of a zero of this function within the interval zL < z < zU.
Since at the two endpoints zL and zU the values of the function are largely away
from zero, we search a zero in a narrower interval than that considered above.
We try zL = 14.79 and compute the values of the Stumpff functions at this point,
as follows

1

1 cos 14:792
CL ¼ C ð14:79Þ ¼ ¼ 0:119144012
14:79
1 1
14:792 sin 14:792
SL ¼ Sð14:79Þ ¼ 3 ¼ 0:0789955175
14:792
The correspondent value of y is
AðzL SL 1Þ
yL ¼ r d þ r a þ 1 ¼ 1:48244650 108 þ 2:19447080 108
CL2
14:79 0:0789955175 1
8:65175147 107 1 ¼ 3:25496359 108 km
0:1191440122
Since yL > 0, we can compute

12 12
yL 3:25496359 108 1
xL ¼ ¼ ¼ 52268:1295 km2
CL 0:119144012
At x = xL and within the interval zL < z < zU, we seek a zero of the following
function
1
x3L SL þ A y2L
f ðDt; xL Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, xL = 52268.1295, SL = 0.0789955175,

1
A = −8.65175147 107, yL = 3.25496359 108, and lS =2 = 3.64297186 105
into the preceding equation, we find
f2 f ðDt; xL Þ ¼ 18276:5778 s
Now, we try zU = 14.80 and compute the values of the Stumpff functions, as
follows

1

1 cos 14:802
CU ¼ C ð14:80Þ ¼ ¼ 0:119006602
14:80

1 1
14:802 sin 14:802
SU ¼ Sð14:80Þ ¼ 3 ¼ 0:0789556894
14:802
AðzU SU 1Þ
yU ¼ rd þ ra þ 1 ¼ 1:48244650 108 þ 2:19447080 108
2
CU
14:80 0:0789556894 1
8:65175147 107 1 ¼ 3:25421722 108 km
0:1190066022
Since yU > 0, we can compute

12 12
yU 3:25421722 108 1
xU ¼ ¼ ¼ 52292:2997 km2
CU 0:119006602
function
1
x3U SU þ A y2U
f ðDt; xU Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, xU = 52,292.2997, SL = 0.0789556894,

1
A = −8.65175147 107, yL = 3.25421722 108, and lS =2 = 3.64297186 105
f1 f ðDt; xU Þ ¼ 9557:29735 s
Since the signs of f(Dt, xL) and f(Dt, xU) are opposite, then the unknown root z of
the time function f(Dt, x) = 0 falls between the lower endpoint zL = 14.79 and the
upper endpoint zU = 14.80. We use Müller’s method of parabolic interpolation to
find iteratively the unknown value of z. This method operates a quadratic inter-
polation and consequently requires three points in the vicinity of the root to be
found. Two of these points are the two endpoints zL z2 and zU z1 found above
and the respective values of the time function. Then, we choose arbitrarily a third
point, z0, between z2 and z1, so that z0 = 14.798.
Then, the function f0, corresponding to z0 = 14.798, is computed as follows

1

1 cos 14:7982
C0 ¼ C ð14:798Þ ¼ ¼ 0:119034075
14:798

1 1
14:7982 sin 14:7982
S0 ¼ Sð14:798Þ ¼ 3 ¼ 0:0789636533
14:7982
Aðz0 S0 1Þ
y0 ¼ rd þ ra þ 1 ¼ 1:48244650 108 þ 2:19447080 108
2
C0
14:798 0:0789636533 1
8:65175147 107 1 ¼ 3:25436646 108 km
0:1190340752
There results y0 > 0; then, we can compute

12 12
y0 3:25436646 108 1
x0 ¼ ¼ ¼ 52287:4639 km2
C0 0:119034075
We seek a zero of the following function

1
x30 S0 þ Ay20
f ðDt; x0 Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, x0 = 52287.4639, S0 = 0.0789636533,

1
A = −8.65175147 107, y0 = 3.25436646 108, and lS =2 = 3.64297186 105
f0 f ðDt; x0 Þ ¼ 3987:01801 s
First iteration. We set the lower endpoint z2, the upper endpoint z1, and the
midpoint z0 as follows
z2 ¼ 14:790 f2 ¼ 18276:5778 s
z0 ¼ 14:798 f0 ¼ 3987:01801 s
z1 ¼ 14:800 f1 ¼ 9557:29735 s
Hence, we compute
h1 ¼ z1 z0 ¼ 14:800 14:798 ¼ 0:002

h2 ¼ z0 z2 ¼ 114:798 14:790 ¼ 0:008
c ¼ h2 =h1 ¼ 0:008=0:002 ¼ 4
interpolates the three points (z2, f2), (z0, f0), and (z1, f1) indicated above
c f1 f0 ð1 þ cÞ þ f2
a¼ ¼ 219019:375
ch21 ð1 þ cÞ
f1 f0 ah21
b¼ ¼ 2:78470163 106
h1
c ¼ f0 ¼ 3987:01801
2c
z ¼ z0 1 ¼ 14:7965681
b ð b2 4acÞ2
to minus, because the value of b is less than zero.
The values of the functions C(z) and S(z) at z = 14.7965681 are

1

1 cos 14:79656812
C0 ¼ C ð14:7965681Þ ¼ ¼ 0:119053746
14:7965681

1 1
14:79656812 sin 14:79656812
S0 ¼ Sð14:7965681Þ ¼ 3 ¼ 0:0789693557
14:79656812
The correspondent value of y (in km) is
Aðz0 S0 1Þ
y0 ¼ r d þ r a þ 1 ¼ 1:48244650 108 þ 2:19447080 108
C02
14:7965681 0:0789693557 1
8:65175147 107 1 ¼ 3:25447332 108
0:1190537462
There results y0 > 0; then, we can compute

12 12
y0 3:25447332 108 1
x0 ¼ ¼ ¼ 52284:0023 km2
C0 0:119053746
We seek a zero of the following function

1
x30 S0 þ A y20
f ðDt; x0 Þ ¼ Dt 1
l2S
After substituting Dt = 2.66976 107, x0 = 52,284.0023, S0 = 0.0789693557,

1
A = −8.65175147 107, y0 = 3.25447332 108, and lS =2 = 3.64297186 105
f ðDt; x0 Þ ¼ 0:0554480480 s
The result found above makes it useless to proceed any further. Therefore, we
accept z0 = 14.7965681 as the solution of the time equation f(Dt, x0) = 0. The
corresponding value of y computed above is 3.25447332 108 km.
The Lagrangian coefficients f, g, and g′ are
y 3:25447332 108
f ¼1 ¼1 ¼ 1:19533947
rd 1:48244650 108
y 3:25447332 108
g0 ¼ 1 ¼ 1 ¼ 0:483033322
ra 2:19447080 108
12 1
y 3:25447332 108 2
g¼A ¼ 8:65175147 10
7
¼ 4:28438585 106
lS 1:32712440 1011
The three components, X 0 d , Y 0 d , and Z 0 d , of the velocity vector, vd, of the

spacecraft in the heliocentric ecliptic reference system XYZ at the time of departure
are
Xa fXd 2:08489513 107 þ 1:19533947 1:04998587 108

X0d ¼ ¼
g 4:28438585 106
¼ 24:4282396 km/s
Ya fYd 2:18416657 108 þ 1:19533947 1:04650719 108
Y 0d ¼ ¼
g 4:28438585 106
¼ 21:7822403 km/s
Za fZd 4:06276688 106 þ 1:19533947 1:12120077 103
Z0d ¼ ¼
g 4:28438585 106
¼ 0:947960059 km/s
The three components, X 0 a , Y 0 a , and Z 0 a , of the velocity vector, va, of the

spacecraft in the heliocentric ecliptic reference system XYZ at the time of arrival are
g0 Xa Xd 0:483033322 2:08489513 107 1:04998587 108

X0a ¼ ¼
g 4:28438585 106
¼ 22:1566993 km/s
g0 Ya Yd 0:483033322 2:18416657 108 1:04998587 108
Y 0a ¼ ¼
g 4:28438585 106
¼ 0:198815989 km/s
g0 Za Zd 0:483033322 4:06276688 106 1:12120077 103
Z0a ¼ ¼
g 4:28438585 106
¼ 0:457785701 km/s
where Xa, Ya, and Za are the three components of the position vector of Mars at the
time of arrival, and Xd, Yd, and Zd are the three components of the position vector of
the Earth at the time of departure in the heliocentric ecliptic reference system XYZ.
These components have been determined above by means of the
NASA-JPL-HORIZONS Web-Interface [19].
The transfer trajectory computed above turns out to be a segment of an ellipse, as
will be shown below. First, the magnitude of the velocity vector vd of the spacecraft
at the moment of departure is

1 h i12
vd ¼ X 0 d þ Y 0 d þ Z 0 d ¼ ð24:4282396Þ2 þ 21:78224032 þ 0:9479600592
2 2 2 2
¼ 32:7429918 km/s
On the other hand, there results

12 1
2lS 2 1:32712440 1011 2
¼ ¼ 42:3137290 km/s
rd 1:48244650 108
Since the values of vd and (2lS/rd)½ computed above are such that
12
2lS
vd \
rd
then the conic transfer trajectory is a segment of an ellipse, as has been shown in
Sect. 1.5.
At the time td of departure from the sphere of influence of the Earth, the position
vector rd of the spacecraft in the heliocentric ecliptic reference system XYZ is the
same as the position vector of the Earth, which results from the JPL data.
At the same time, the velocity vector vd of the spacecraft has the components
X 0 d , Y 0 d , and Z 0 d computed above. For convenience, the components of these
vectors (expressed in km and km/s) are written below.
rd ¼ 1:04998587 108 uX þ 1:04650719 108 uY þ 1:12120077 103 uZ

vd ¼ 24:4282396 uX þ 21:7822403 uY þ 0:947960059 uZ
Likewise, at the time ta of arrival at the sphere of influence of Mars, the com-
ponents of the position vector ra and the velocity vector va of the spacecraft (ex-
pressed in km and km/s) in the heliocentric ecliptic reference system XYZ are
ra ¼ 2:08489513 107 uX 2:18416657 108 uY 4:06276688 106 uZ

va ¼ 22:1566993 uX 0:198815989 uY 0:457785701 uZ
Either of these two sets of six components can be used to compute the six orbital
elements of the transfer trajectory of the spacecraft. By using the first set and
applying the method shown in Sect. 1.9, we obtain
a ¼ 1:84762129 108 km e ¼ 0:20586209 i ¼ 1 :66193782

X ¼ 44 :8880895 x ¼ 19 :9930656 /d ¼ 340 :021866
At the time td of departure from the sphere of influence of the Earth, the
hyperbolic excess velocity vector (in km/s) of the spacecraft is
vd1 ¼ vd vEd
where vEd is the orbital velocity of the Earth with respect to the Sun at t = td.
The components of vd have been computed above. The components of vEd result
from the JPL data. Therefore, we have
vd1 ¼ vd vEd ¼ ð24:4282396 þ 21:5149121Þ uX þ ð21:7822403 20:9988238Þ uY

þ ð0:947960059 þ 0:000948089222Þ uZ
¼ 2:9133275 uX þ 0:7834165 uY þ 0:94890815 uZ

h i12
vd1 ¼ ð2:9133275Þ2 þ 0:78341652 þ 0:948908152 ¼ 3:1625378 km/s
Likewise, at the time ta of arrival at the sphere of influence of Mars, the

hyperbolic excess velocity vector (in km/s) of the spacecraft is
va1 ¼ va vMa
where vMa is the orbital velocity of Mars with respect to the Sun at t = ta.
The components of va have been computed above. The components of vMa result
from the JPL data. Therefore, we have
va1 ¼ va vMa ¼ ð22:1566993 25:0368939ÞuX þ ð0:198815989 þ 0:220527950ÞuY

þ ð0:457785701 þ 0:620131431ÞuZ
¼ 2:8801946 uX þ 0:021711961 uY þ 0:16234573 uZ

h i12
va1 ¼ ð2:8801946Þ2 þ 0:0217119612 þ 0:162345732 ¼ 2:8848481 km/s
In case of the spacecraft being placed, before the time of departure, in a circular
parking orbit of radius 6378.137 + 200 = 6578.137 km around the Earth, we want
to compute the velocity increment necessary to insert the spacecraft into the
departure hyperbola.
As has been shown in Sect. 5.6, the velocity of a spacecraft at the periapsis of
the hyperbola of departure from a planet is
12
2lPd
vP ¼ þ vd1
2
rP
where lPd is the gravitational parameter, rP is the radius of the periapsis of the
planet of departure, and vd∞ is the magnitude of the hyperbolic excess velocity
vector.
In the present case, the planet of departure is the Earth, whose gravitational
parameter is lE = 398600.44 km3/s2, rP = 6578.137 km, and vd∞ has been found
above to be equal to 3.1625378 km/s.
By inserting these values into the preceding expression, there results
12
2 398600:44
vP ¼ þ 3:1625378 ¼ 11:453869 km/s
2
6578:137
On the other hand, the velocity of the spacecraft in its circular parking orbit is
12 1
lE 398600:44 2
vC ¼ ¼ ¼ 7:7842617 km/s
rP 6578:137
Therefore, the velocity increment to be given to the spacecraft is
Dv ¼ vP vC ¼ 11:453869 7:7842617 ¼ 3:6696073 km=s
The eccentricity of the departure hyperbola results from the following expression
derived in Sect. 5.6
rP v2d1
e¼ þ1
lPd
6578:137 3:16253782
e¼ þ 1 ¼ 1:1649576
398600:44
Finally, we compute the variation of velocity necessary in order for the space-
craft to be captured by the gravitational field of Mars in an elliptic orbit, such that
the period of revolution be T = 48 h = 172800 s and the altitude of the periapsis be
h = 300 km.
In accordance with NASA [20], the mean radius and the gravitational parameter
of Mars are assumed to be rM = 3389.5 km and lM = 42830 km3/s2.
The radius of the periapsis of the hyperbola of arrival is the same as that of the
capture ellipse, that is,
rP ¼ rM þ h ¼ 3389:5 þ 300 ¼ 3689:5 km
As has been shown in Sect. 5.7, the velocity of the spacecraft at the periapsis of
the hyperbola of arrival is
12
2lPa
vP ¼ þ va1
2
rP
where lPa is the gravitational parameter of the planet of arrival, rP is the radius of
periapsis of the hyperbola of arrival, and va∞ is the magnitude of the hyperbolic
excess velocity vector.
In the present case, the planet of arrival is Mars, whose gravitational parameter is
lM = 42830 km3/s2, and va∞ has been found above to be equal to 2.8848481 km/s.
By inserting these values into the preceding expression, there results
12
2 42830
vP ¼ þ 2:88484812 ¼ 5:6160116 km/s
3689:5
The major semi-axis, aE, of the capture ellipse around Mars results from the third
law of Kepler
3 12
a
T ¼ 2p E
lM
The preceding equation, solved for aE, yields

" #13 " #13
T 2 172800 2
aE ¼ l M ¼ 42830 ¼ 31878:064 km
2p 2p
The eccentricity, eE, of the capture ellipse results from
rP 3689:5
eE ¼ 1 ¼1 ¼ 0:88426211
aE 31878:064
The velocity, vPE, of the spacecraft at the periapsis of the capture ellipse results
from the vis-viva integral
12 12
2 1 2 1
vPE ¼ lM ¼ 42830 ¼ 4:6769307 km/s
r P aE 3689:5 31878:064
The variation of velocity necessary in order for the spacecraft to be captured by

the gravitational field of Mars in the desired elliptic orbit is
Dv ¼ vP vPE ¼ 5:6160116 4:6769307 ¼ 0:9390809 km=s
Since vP is greater than vPE, then the variation of velocity necessary for this
purpose is a deceleration, which may be performed by either a motor rocket burn in
the direction opposite to that of the orbital velocity vector, or by aerodynamic
means, as will be shown in the following paragraph.
The eccentricity, e, of the hyperbola of arrival has been found in Sect. 5.7 to be
rP v2a1
e¼ þ1
lPa
where lPa is the gravitational parameter of the planet of arrival.

In the present case, the planet of arrival is Mars, whose gravitational parameter is
lM = 42830 km3/s2, and the magnitude, va∞, of the hyperbolic excess velocity
vector has been found above to be 2.8848481 km/s. By inserting these values into
the preceding expression, there results
3689:5 2:88484812
e¼ þ 1 ¼ 1:7189112
42830
5.12 The Aerodynamic Assist
Aerodynamic assist manoeuvres are meant to reduce the propellant required to

insert a spacecraft into a desired final orbit around a target planet whose atmosphere
is sufficiently dense for this purpose. There are two principal techniques which
operate aerodynamic deceleration. The first of them, called aerobraking and illus-
trated in the following figure, due to the courtesy of NASA [21], is a gradual
process, which reduces the velocity of a spacecraft by means of the drag acting on
the main body and the solar arrays of the spacecraft passing through the upper
atmosphere of a target planet and avoids an increase of mass due to a heat shield.
The second technique, called aerocapture, is a very rapid process, which requires
a heavy heat shield to withstand the heat produced by the atmospheric friction.
These techniques are briefly described below.
As has been shown in Sects. 3.19 et seq., the atmospheric drag slows down a
spacecraft passing through a planet which has an atmosphere. In case of a planet
having an atmosphere of sufficient density, this effect can be used intentionally to
either lower the altitude of an orbit or decrease the velocity of a spacecraft before
landing.
For this purpose, a spacecraft approaching a planet along an arrival hyperbola
performs a propulsive manoeuvre to be captured into an elliptic orbit whose peri-
apsis is inside the atmosphere of the target planet. The aerobraking technique uses
the atmospheric drag to reduce gradually the velocity of the spacecraft at the
periapsis of the capture ellipse in order to lower the apapsis of this ellipse.
5.12 The Aerodynamic Assist 839
Hundreds of passes of the spacecraft through the atmosphere of the planet are
necessary to reduce the altitude of the apapsis to the desired value. When this value
is reached, another propulsive burn is performed at the apapsis in order to raise the
periapsis above the atmosphere and circularise the orbit. The time required by this
process is kept as small as possible, in order to reduce the heat generated and avoid
the installation of a heat shield. So far, this technique has been used in four space
missions: Magellan, Mars Global Surveyor, Mars Odyssey, and Mars
Reconnaissance Orbiter [22].
The aerocapture technique is illustrated in the following figure, which is also due
to the courtesy of NASA [21].
As is the case with aerobraking, aerocapture uses the atmosphere of a target

planet to slow an approaching spacecraft and insert it into an elliptic orbit around
the planet. However, it differs from aerobraking because the spacecraft is inserted
from the hyperbola of arrival to the elliptic orbit as soon as it arrives at the target
planet. In other words, aerocapture occurs in a single atmospheric pass instead of
hundreds of passes, and therefore, the spacecraft is inserted immediately into the
elliptic orbit. The variation of velocity is due for the greatest part (over 95%) to
aerodynamic drag.
This manoeuvre begins at the moment in which the spacecraft enters the
atmosphere of the target planet from its hyperbola of arrival. A quick, nearly
inexpensive (from the point of view of propellant consumption) motor firing causes
the spacecraft to enter a dense layer of the atmosphere. The friction between the
spacecraft and the atmosphere decelerates the spacecraft and inserts it into an
elliptic orbit after the first pass through the periapsis. Then, the spacecraft exits the
atmosphere and reaches the apapsis of the same orbit, at which point a second
motor firing is performed in order to raise the periapsis and circularise the orbit. An
aerocapture manoeuvre requires one of two possible configurations to protect the

spacecraft from the heat developed by atmospheric friction. For this purpose, the
spacecraft may be enveloped by a structure having heat shields applied to its
surfaces. Rigid shells of this type may be, in turn, either blunt bodies (such as those
used in the Apollo Command Module, or more recently in the Mars Exploration
Mission Spirit and Opportunity, which landed on the surface of Mars in January
2004) or slender bodies, that is, elongated capsules having a hard shell surrounding
the spacecraft. Otherwise, the spacecraft may deploy an aerocapture device, such as
an inflatable heat shield or an inflatable trailing ballute, the latter being a combi-
nation of a balloon and a parachute made of thin and durable material towed behind
the spacecraft after deployment in space [23].
The following figure, due to the courtesy of NASA [24], shows a toroidal
ballute, which is much larger than the spacecraft and is towed like a parachute to
slow the vehicle down. A ballute is made of a lightweight, thin-film material and
can be detached after the completion of an aerocapture manoeuvre.
The principal advantage offered by aerobraking or aerocapture manoeuvres

resides in propellant mass savings. Such savings, in turn, result in smaller launch
vehicles or increased payloads.
5.13 Trajectories of Vehicles Propelled by Solar Radiation Pressure 841
5.13 Trajectories of Vehicles Propelled by Solar Radiation

Pressure
As has been shown in Sect. 3.16, the solar radiation pressure is due to the radiation
emitted by the Sun. This radiation imparts a momentum on the surface of a body
exposed to it, and this transfer of momentum results in a force acting upon that
body. The radiation incident on a body is partially absorbed, partially reflected
specularly, and partially reflected diffusively by the illuminated body. In the ideal
case of a body having a perfectly specular surface, all the radiation impinging on
that body can be used to accelerate it.
It has also been shown that the energy, E, of a quantum of electromagnetic
radiation is proportional to the frequency, f, of the radiation itself, as follows
E ¼ hf
where h = 6.62606957 10−34 J s is the Planck constant. Einstein derived a

relation between the energy, E, and the momentum, p, of a quantum (photon) of
light energy, as follows
E ¼ cp
where c = 299792458 m/s is the speed of light in vacuo.

hf
p¼
c
The preceding expression, integrated over all the frequencies and further inte-
grated over the number of photons per unit area, yields the force per unit area. At
the distance, rE, of one Astronomical Unit (or 149597870700 m) from the Sun, this
force per unit area is
S0
P¼
c
where S0 is the total solar irradiance (TSI). Consequently, the total solar irradiance
is the amount of radiant energy emitted by the Sun over all the frequencies (or
wavelengths) which falls each second on the unit area (1 m2) exposed normally to
the Sun rays at the mean Sun–Earth distance in the absence of the Earth atmo-
sphere. Since the Earth orbit around the Sun is not circular but slightly elliptical,
then the intensity of the solar radiation received in a point outside the Earth
atmosphere varies with the square of the Earth–Sun distance.
On average, the total solar irradiance (also called solar constant) is

Ssc = 1360.8 Watts per square metre (W/m2), according to Kopp and Lean [25].
The value of S0 varies by about ±3.4% with respect to this value, for the reason
indicated above.
By substituting S0 Ssc = 1360.8 W/m2 into the preceding equation, there
results
S0 1360:8
P¼ ¼ 4:54 106 N/m2
c 299792458
which is an approximate value of the solar radiation pressure acting on a perfect

reflector placed at a distance of 1 AU from the Sun.
Therefore, in order for the magnitude of the corresponding force to be sufficient
to propel a spacecraft, a large reflecting area, A, called solar sail, is necessary.
Spacecraft propelled by means of solar sails should have thin-film (2–5 lm), large
(e.g. the Sunjammer mission, cancelled in October 2014, was designed to have a
solar sail measuring 38 38 m2), low-density (7–12 g/m2), reflective surfaces
made usually of aluminium-coated Kapton® or Mylar®, and these surfaces should
be easily deployable in space.
The total solar irradiance is related to the luminosity of the Sun, LS, as follows
LS
S0 ¼
4prE2
At any other distance, r, than rE = 1 AU, the amount of radiant energy, S, emitted
by the Sun over all the frequencies, which falls each second on the unit area (1 m2)
exposed normally to the Sun rays is

r 2
E
S ¼ S0
r
The definition of S given above implies that
1 DE
S¼
A Dt
Therefore, a variation DE of radiant energy carries a variation of momentum
DE ASDt
Dp ¼ ¼
c c
Newton’s second law of dynamics states that the force corresponding to this
variation of momentum is
Dp AS
F¼ ¼
Dt c
Therefore, the solar radiation pressure at a distance r from the Sun is
F S
P¼ ¼
A c
in accordance with the equation derived above for r = rE.

There are at present three principal types of solar sails. They are shown in the
following figure, which is due to the courtesy of NASA [26].
The first type, called square sail, has four deployable spars cantilevered from a
central hub. The sail attitude is controlled on three axes by attaching articulated
reflecting vanes of variable pitch at the tips of the spars, for the purpose of gen-
erating torques. The sails of this type have the advantage of being testable on the
ground, at least for those of small size, but the disadvantage of having high masses,
due to their supporting structures, which are usually booms or deployable trusses.
The second type is a disc-shaped sail, where a thin membrane is held in tension
by spin. The sails of this type are deployed and flattened by centrifugal forces only.
They are lighter than those of the first type, but have problems of attitude control of
the spinning sail membrane.
The third type, called heliogyro, is a structure which has long (of the order of
magnitude of kilometres) reflective membrane strips, called usually blades, which
extend from a central spinning hub and produce thrust from solar radiation pressure.
The centrifugal force extends the blades outward and keeps them stiff. A heliogyro
is also lighter than a sail of the first type. Its blades are stowed and deployed from
reels. Attitude control is obtained by pitching the blades at the root to change their
orientation with respect to the Sun [27].
Solar sails can be oriented with respect to the Sun, and therefore, the acceleration
of a spacecraft depends on the sail attitude, as will be shown below.
With reference to the preceding figure, a ray of light coming from the Sun carries
momentum incident on the reflector. According Newton’s second law, a change of
momentum, due to the reflected ray, causes an applied force. Therefore, according
Newton’s third law, a reaction force acts upon the reflector.
Let ui, ur, and un be the unit vectors of, respectively, the incident momentum,
the reflected momentum, and the normal to the reflecting surface. The force exerted
on this surface, due to the incident photons from the direction ui, is
SAðui un Þ SA cos a
Fi ¼ ui ¼ ui
c c
where a is the angle, measured with respect to the normal un to the reflector,
through which the incoming ray of light is reflected.
Likewise, the force exerted on the same surface, due to the photons reflected in
the direction ur, is
SAður un Þ SA cos a
Fr ¼ ur ¼ ur
c c
and has the same magnitude as that of Fi, but is directed along the unit vector ur of
the reflected ray. Since
ui ur ¼ 2ðui un Þun ¼ ð2 cos aÞun

is an identity [28], then the total reaction force exerted on the reflector is
2 SA cos2 a
F¼ un
c
which is directed along the normal unit vector un. It is convenient to decompose the
force vector F, or the unit vector un, in its radial (in the direction of ur), transverse
(in the direction of uw), and bi-normal (in the direction of uh) components, which
are shown in the following figure.
The unit vector ur = r/r has been defined above. The unit vector uw lies in the
instantaneous plane r, v of motion of the spacecraft, is perpendicular to r, and points
in the direction of motion (uw v > 0, where v is the velocity vector of the spacecraft
with respect to the Sun). The unit vector uh = h/h (where h = r v is the moment of
momentum vector per unit mass of the spacecraft) is perpendicular to the instan-
taneous plane of motion of the spacecraft.
The three unit vectors ur, uw, and uh form a right-handed system of reference.
The angle d shown in the preceding figure, called the sail clock angle, is the angle
which the projection of the unit vector un onto the plane uw, uh forms with the
direction of uh. Therefore, the unit vector un can be expressed as follows
un ¼ ðcos aÞ ur þ ðsin a sin dÞ uw þ ðsin a cos dÞuh
Since the force vector F is directed along un, then its components Fr, Fw, and Fh
are proportional to, respectively, cos a, sin a sin d, and sin a cos d.
In the particular case of a spacecraft which travels constantly in the ecliptic plane
(d = p/2), the component Fh along uh is equal to zero, and consequently, the force
vector F can be expressed as follows
F ¼ F r ur þ Fw uw ¼ ðF cos aÞur þ ðF sin aÞuw
Since the force vector F has been shown to be
2 SA cos2 a
F¼ un
c
then Fr is proportional to cos3 a, and Fw is proportional to cos2a sin a.

There is an optimum value, a*, of the sail cone angle a, such that the transverse
component, Fw, of the force vector F has the maximum absolute value. In order to
determine this optimum value, we differentiate the quantity cos2a sin a (to which
the transverse component Fw is proportional) with respect to a, and set the
derivative equal to zero. This yields
2 cos a sin2 a þ cos3 a ¼ 0
The preceding cubic equation has the root cos a = 0. Discarding this root, we
search the two further roots resulting from
2 sin2 a þ cos2 a ¼ 1 3 sin2 a ¼ 0
The preceding quadratic equation, solved for sin a, yields

12
1
sin a ¼
3
which in turn yields the optimum value of the sail cone angle:
a 35 :26
In other words, the value of the sail cone angle, a, should be kept constant and as
close as possible to ±35°.26 during the space travel, in order to obtain the maxi-
mum absolute value of the transverse component, Fw, of the force F, which is due
to the solar radiation pressure. The choice of the sign (plus or minus) in front of
35°.26 (i.e. the choice of the orientation of the solar sail with respect to the incident
ray coming from the Sun) depends on whether a spacecraft is desired to travel along
an outward spiral (to go from the Earth to Mars) or along an inward spiral (to go
from the Earth to Venus), as shown on, respectively, the left-hand and the
right-hand sides of the following figure.
The choice of sign for the sail cone angle a has no effect on the radial component
Fr, which is proportional to cos a, and therefore is always directed outward, that is,
away from the Sun, as shown in the following figure.
As has been shown by Prodger [29], space vehicles propelled by solar sails have
some indicators of performance, which are briefly described below.
The first of them, denoted by r, is the sail loading, which is the mass per unit
area of the overall spacecraft
m
r¼
A
and is usually expressed in grammes per square metre (g/m2).

The second indicator, denoted by ac, is the characteristic acceleration, which is
the acceleration, due to solar radiation pressure, acting on a solar sail facing the Sun
(a = 0) at a distance of one astronomical unit.
Remembering the equation
2 SA cos2 a
F ¼ F un ¼ un
c
and setting F = ma and A = m/r, there results
2S cos2 a
a¼
rc
which accounts for the definition of characteristic acceleration, ac, given above.
This definition does not take account of the sail efficiency, denoted by η, whose
value depends on the type of material used for the reflecting membrane, on the
manufacturing accuracy of this membrane, and also on creases which may occur
during its deployment. Since aluminium has a reflectivity ranging from 86.8% (at a
wavelength of 0.8 lm) to 92% (at a wavelength of 0.4 lm) [30], then a value 0.85
may be taken for the efficiency of a membrane coated with aluminium [29]. Taking
account of the sail efficiency, η, the characteristic acceleration is defined as follows
2S g
ac ¼
rc
and is usually expressed in mm/s2.

The third indicator of performance, denoted by b, is the sail lightness number,
which is defined as the ratio of the solar radiation pressure acceleration to the solar
gravitational acceleration.
Since the solar radiation pressure acceleration, at a distance r from the Sun, is
2S0 g
rE 2
aSRP ¼
rc r
and the solar gravitational acceleration, at the same distance, is

lS
aG ¼
r2
where lS = 1.327 1020 m3/s2 [2] is the gravitational parameter of the Sun, then
the sail lightness number is
aSRP 2S0 grE2

b¼ ¼
aG rc lS
After substituting S0 Ssc = 1360.8 W/m2, η = 0.85, c = 299792458 m/s, rE =

149597870700 m, and lS = 1.327 1020 m3/s2 into the preceding equation, we
find
0:0013
b
r
The sail lightness number b is dimensionless. The numerator of the fraction

written above (r* = 0.0013 kg/m2 = 1.3 g/m2) is called the critical sail loading.
A spacecraft, if it were equipped with a solar sail capable of generating a sail
loading equal to 1.3 g/m2, would remain suspended in space, because the solar
radiation pressure acceleration acting upon it would be equal to its solar gravita-
tional acceleration (b = 1).
The following part of this paragraph deals with the trajectory of a spacecraft
propelled by solar radiation pressure. As has been shown above, the radial (ar),
transverse (aw), and bi-normal (ah) components of the acceleration vector a due to
solar radiation pressure are expressible as follows
2 3 2 3
ar
cos a
4 aw 5 ¼ r 2
ac cos2 a 4 sin a sin d 5
E
r
ah sin a cos d
where ac is the characteristic acceleration. These components are to be introduced

into the Lagrange planetary equations in Gaussian form (see Sect. 3.2):
2a2 h
p i
a0 ¼ ðe sin /Þar þ aw
h r
1
e0 ¼ ½p sin /
ar þ ½ðp þ r Þcos / þ re
aw
h

r cosðx þ /Þ
i0 ¼ ah
h

r sinðx þ /Þ
X0 ¼ aw
h sin i

p cos / ðp þ r Þsin / r sinðx þ /Þcos i
x0 ¼ ar þ aw ah
he he h sin i
h 1
/0 ¼ 2þ ½p cos /
ar ½ðp þ r Þsin /
aw
r eh
where the prime sign (′) denotes first derivatives with respect to time, / is the true
anomaly of the spacecraft in the instantaneous osculating orbit, p = a(1 − e2) is the
semi-latus rectum, h = (lSp)½ is the magnitude of the moment of momentum per
unit mass of the spacecraft, and lS is the gravitational parameter of the Sun. The
preceding differential equations, integrated numerically with the initial conditions,
give the six classical elements (see Sect. 1.9) of the osculating orbit.
The third of these equations shows that the rate of change (i′) of the inclination
angle (i) of the instantaneous osculating orbit with respect to the ecliptic plane
depends on cos(x + /) and also on the bi-normal acceleration ah, whose value in
turn is proportional (see above) to cos2a sin a cos d. The optimum value of the cone
angle a has been shown to be a* = ±35°.26. When it is desired to use the solar
radiation pressure for the purpose of changing the value of the angle i in the most
efficient manner, then it is necessary to impose the following condition
d
cosðx þ /Þcos2 a sin a cos d ¼ 0
dd
Since sin d is equal to zero when d is equal to either 0 or p, then Wie [31] has
proposed the following steering law to obtain the maximum rate of change of the
inclination angle:
d¼0 for cos ðx þ /Þ 0

d¼p for cos ðx þ /Þ\0
As an option, the differential equations governing the motion of a spacecraft

propelled by solar radiation pressure may be written in heliocentric ecliptic rect-
angular co-ordinates X, Y, and Z. In this case, it is necessary to take account not only
of the solar radiation pressure but also of the gravitational acceleration, as follows
lS X lS Y lS Z
X 00 ¼ þ aX Y 00 ¼ þ aY Z 00 ¼ þ aZ
r3 r3 r3
where r = (X2 + Y2 + Z2)½. The three components aX, aY, and aZ of the acceleration
vector a, due only to the solar radiation pressure, result from
2 3 2 3
aX ar
4 aY 5 ¼ R3 ðXÞR1 ðiÞR3 ðx /Þ4 aw 5
aZ ah
where ar, aw, ah are the three components of a which have been expressed above,
and R1 and R3 are two of the following three elementary rotation matrices
2 3 2 3
1 0 0 cos h 0 sin h
6 7 6 7
R1 4 0 cos h sin h 5 R2 4 0 1 0 5
0 sin h cos h sin h 0 cos h
2 3
cos h sin h 0
6 7
R3 4 sin h cos h 0 5
0 0 1
(h being any angle), as has been shown in Sect. 1.9.

As another option, the same differential equations may be written in heliocentric

spherical co-ordinates r, w, and u, which are, respectively, the distance, the ecliptic
longitude and the ecliptic latitude of the point occupied at a time t by the spacecraft.
The ecliptic longitude (0 w 2p) is measured along the ecliptic, increasing to
the east, with the zero point at the vernal equinox of some standard epoch (e.g. of
J2000.0). The ecliptic latitude (−p/2 u p/2) is measured along the local
ecliptic meridian from the ecliptic plane to the ecliptic poles. Let ar, aw, and au be
the three components of the acceleration vector a due to the solar radiation pressure.
These components are expressible as follows
2 3 2 3
ar aX
4 aw 5 ¼ R2 ðuÞR3 ðwÞ4 aY 5
au aZ
By using the heliocentric spherical co-ordinates, the differential equations gov-

erning the motion of a spacecraft propelled by solar radiation pressure may be
written as follows
lS
r 2
r 00 ru02 rw02 cos2 u ¼
E
þ a c cos3 a
r2 r

r 2
r w00 cos u þ 2r 0 w0 cos u 2rw0 u0 sin u ¼ ac
E
cos2 a sin a sin d
r

r 2
r u00 þ 2r 0 u0 þ rw02 sin u cos u ¼ ac
E
cos2 a sin a cos d
r
The equations of motion written above can be applied to the particular case of a
logarithmic (also called equiangular) spiral trajectory laying in the ecliptic plane.
For this trajectory, proposed by Bacon in 1959 [32], the angle between the tangent
and the radial line at any point is constant. Therefore, a logarithmic spiral lends
itself to a simple steering law with a constant value of the sail cone angle a.
In case of u = 0 and d = p/2, the preceding equations of motion reduce to
lS
r 2
r 00 r w02 ¼
E
þ a c cos3 a
r2 r

r 2
r w00 þ 2r 0 w0 ¼ ac
E
cos2 a sin a
r
Remembering the definition (b = aSRP/aG) given above for the sail lightness
number, the two preceding equations can be written as follows
lS l
r 00 r w02 ¼ 2
þ b 2S cos3 a
r r
l
r w00 þ 2r 0 w0 ¼ b 2S cos2 a sin a
r
A particular solution of the two differential equations written above is sought in

the following form
r ¼ r0 expðw tan cÞ
where c is the constant flight path angle formed by the velocity vector v of the
spacecraft with the transverse unit vector uw. The preceding equation is the equation
of a logarithmic spiral in polar co-ordinates r and w.
The derivatives r′ and r″ of the preceding function r with respect to time are
r 0 ¼ ½r0 expðw tan cÞ

w0 tan c ¼ r w0 tan c
0
r 00 ¼ ðr w0 Þ tan c ¼ ðr 0 w0 þ r w00 Þtan c ¼ r w02 tan c þ w00 tan c
After substituting these derivatives in the two equations of motion written above,
there results
l l
r w02 tan c þ w00 tan c r w02 ¼ 2S þ b 2S cos3 a
r r
l
r w00 þ 2w02 tan c ¼ b 2S cos2 a sin a
r
The second of these equations is solved for w″ and the resulting expression is
substituted into the first equation. This yields after simplification

r 3 w02 ¼ lS 1 bðcos a tan c sin aÞcos2 a cos2 c
Since the quantities on the right-hand side of the preceding equation are con-
stant, then the product r3w′2 is also constant, in case of a spacecraft which moves
with respect to a principal attracting body along a logarithmic spiral trajectory.
Taking account of the preceding equation and also of r′ = rw′ tan c, the radial
component (vr), the transverse component (vw), and the magnitude (v) of the
instantaneous velocity vector v of a spacecraft moving along a logarithmic spiral
trajectory are

l 12 1
vr ¼ r 0 ¼ S
1 bðcos a tan c sin aÞcos2 a 2 sin c
r

l 12 1
vw ¼ rw0 ¼ S 1 bðcos a tan c sin aÞcos2 a 2 cos c
r

12
l 12 1
v ¼ v2r þ v2w ¼ S 1 bðcos a tan c sin aÞcos2 a 2
r
Since the value of [1 − b(cos a − tan c sin a)cos2a]½ is always less than unity,
then the velocity v is less than vcircular, where vcircular = (lS/r)½ is the velocity of a
spacecraft which moves with respect to the Sun along a circular orbit of radius r. By
the way, a transfer between two coplanar concentric circular orbits of different radii
requires: (a) an impulse at the moment of leaving the circular orbit of departure, in
order to insert the spacecraft into the desired logarithmic spiral, and (b) another
impulse at the moment of arrival, in order to insert the spacecraft into the circular
orbit of destination. These impulses imply an expenditure of propellant.
London [33], Wie [31], and McInnes [34] have shown that there is an implicit
formula which relates the sail cone angle a to the flight path angle c, as follows
sin c cos c b cos2 a sin a

¼
2 sin2 c 1 b cos3 a
In practical cases, the value of the angle c is little different from zero, and
therefore the preceding formula can be approximated as follows
2b cos2 a sin a
tan c
1 b cos3 a
Since
Zr Zt
1 1
r dr ¼
2 2lS b sin a cos2 a tan c 2 dt
r0 t0
and
Zr
1 2
3 3

r 2 dr ¼ r 2 r02
3
r0
then the transfer time t − t0 necessary for a spacecraft which moves along a log-
arithmic spiral trajectory to go from a point P0 (placed at a distance r0 from the Sun)
to another point P (placed at a distance r from the Sun) can be computed as follows
1
3 3
2 cot c
12
t t0 ¼ r r0
2 2
3 lS b sin a cos2 a
By introducing the approximate formula tan c (2b sin a cos2a)/(1 − b cos3a)

into the preceding equation, there results
1
3 3
1 b cos3 a 12
t t0 ¼ r 2 r02
3 lS b2 sin2 a cos4 a
The transfer time t − t0 is proportional to the factor

12
1 b cos3 a
lS b2 sin2 a cos4 a
which depends only on the sail lightness number b, because the value of the sail
cone angle a is constant, for the reason indicated above. In order to find the value of
b corresponding to the minimum transfer time, the preceding factor is differentiated,
and its derivative is set equal to zero.
This condition leads to the following expression of the sail lightness number as a
function of the sail cone angle
2 4 tan2 a
b¼
ð2 tan2 aÞcos a
So far, the solar sail technology has been tested in space by NASA’s
NanoSail-D, launched on 19 November 2010 to orbit around the Earth, and by
JAXA’s IKAROS (Interplanetary Kite-craft Accelerated by Radiation Of the Sun),
launched on 21 May 2010 and placed into an interplanetary trajectory towards
Venus. On 8 December 2010, IKAROS passed by Venus at a distance of about
80800 km, completing successfully the planned mission. IKAROS had a square
solar sail measuring 14 14 m2. The force measured by JAXA scientists on the
solar sail was 1.12 10−3 N 1 gf [35].
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uk/*dma0rcj/Psling/sling.pdf
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Chapter 6
Numerical Integration of the Equations
of Motion
As has been shown in Sect. 1.14, the n-body problem (which consists in deter-
mining the motion of an isolated set of n bodies, having masses m1, m2, …, mn, and
attracting one another with Newtonian gravitational forces) cannot be solved ana-
lytically in the general case, when n is greater than two.
For this purpose, it is necessary to use methods which integrate numerically the
differential equations of motion with given initial conditions.
A particular class of these methods is based on the concept of finite differences.
Let us consider a first-order differential equation
x0 ðtÞ ¼ f ½xðtÞ; t
where f[x(t), t] is the integrand function, and x(t) and x′(t) are, respectively, the
unknown function and its first derivative with respect to the independent variable t.
Let a t b be the domain of definition of the preceding differential equation,
and let x0 = x(t0) be the initial condition associated with it.
These methods approximate the first derivative x′(t) of the unknown function x(t)
to the corresponding difference quotient, that is,
x ð t þ hÞ x ð t Þ
x0 ð t Þ
h
where the increment h is much less than unity. By so doing, the given differential
equation is approximated to the following finite-difference equation
xðt þ hÞ ¼ xðtÞ þ h f ½xðtÞ; t
in which the value of x(t + h) at the future time t + h is computed by using the
values of x(t) and the integrand function f[x(t), t] at the present time t.

https://doi.org/10.1007/978-3-319-62220-0_6
858 6 Numerical Integration of the Equations of Motion
In the general case, for a differential equation in vector form
x0 ðtÞ ¼ f ½xðtÞ; t
with the initial condition x0 = x(t0), defined in a domain a t b, a numerical

method based on finite differences makes it possible to calculate the
dependent-variable vector x(t) at a sequence of values t1, t2, …, tm greater than the
initial value t0. In other words, the evaluation of the first-derivative vector x′(t) is
performed only at discrete points within the domain of the definition of the given
differential equation. By considering a first-order differential equation, the follow-
ing discussion is not deprived of generality, because any differential equation
xðnÞ ðtÞ ¼ f ½xðtÞ; t
of order n > 1, which can be expressed in explicit form, that is, solved for its
highest derivative, may be written as a system of n first-order differential equations
x′(t) = f[x(t), t] by an appropriate choice of the dependent-variable vector x(t). For
example, the following second-order differential equation
x00 ðtÞ þ xðtÞ ¼ 0
with the associated initial condition
x ð 0Þ ¼ 0
x 0 ð 0Þ ¼ 1
may be written as a system of two first-order differential equations as follows
x0 1 ðtÞ ¼ x2 ðtÞ
x0 2 ðtÞ ¼ x1 ðtÞ
with the initial condition
x 1 ð 0Þ ¼ 0
x 2 ð 0Þ ¼ 1
Therefore, the solution of a differential equation of the nth order x(n)(t) = f[x(t), t] with
n associated initial conditions xðt0 Þ ¼ x0 ; x0 ðt0 Þ ¼ x0 0 ; . . .; xðn1Þ ðt0 Þ ¼ xðn1Þ 0 is
amenable to the solution of a system of n first-order differential equations in vector
form x′(t) = f[x(t), t] with the n initial conditions x(t0) = x0.
All numerical methods based on finite differences compute approximate solu-
tions of a differential equation in vector form x(′t) = f[x(t), t], with an associated
initial condition x(t0) = x0, to be integrated over an interval a t b, at certain
points, called mesh points, of the integration interval.
Such approximate solutions are sets of tabulated values of the

dependent-variable vector x(t), that is, approximations to the true solution computed
at some progressive values (ti+1, ti+2, …, ti+m) of the independent variable t.
Starting from an approximate solution x(ti) at the point ti contained in the
interval a t b, a method based on finite differences computes another
approximate solution x(ti+1) at the point ti+1 also contained in a t b. The
quantity hi = ti+1 − ti is called the step size and is in general variable from one point
ti to another. The numerical methods for obtaining such approximate solutions
require the expression of the differential equation in a discrete form and the use of
truncated Taylor series and finite-difference formulae. The category of
finite-difference methods comprises two sub-categories. The methods of the first,
called single-step methods, compute approximate solutions at ti+1 without regard to
the values of the unknown function or its derivative at mesh points (ti−1, ti−2, …, ti
−m) preceding the current mesh point (ti). The most popular among single-step
methods, those of the Runge–Kutta type, evaluate the integrand function at some
intermediate points (also called support abscissae) placed between the current mesh
point (ti) and the next mesh point (ti+1); then, a weight is assigned to each of these
values; finally, the value of x(t) at the current mesh point (ti) and a weighted average
of the values of the integrand function at the intermediate points are used in a
formula which computes x(t) at the next mesh point (ti+1). The methods of the
second sub-category, called multi-step methods, require starting values from several
mesh points (ti−1, ti−2, …, ti−m) preceding the current mesh point (ti).
The numerical methods which will be considered in this paragraphs are limited
to those which are most frequently used in astrodynamics. They are the Runge–
Kutta methods and the Bulirsch–Stoer methods, which belong to the first
sub-category, and the Adams, Störmer–Cowell, and Gauss–Jackson methods,
which belong to the second sub-category. The 3 1 position and velocity vectors
are indicated here with, respectively, x and y x′, where the prime sign indicates
the first derivative with respect to time. The 3 1 acceleration vector is indicated
with either a or y′ or x″. The 6 1 state vector is indicated with z.
The step-by-step methods used to solve differential equations are based on a
mathematical theory which leads to an approximate solution for each integration
step. This is principally because the exact solution is expressed by using infinite
difference tables or infinite series. On the contrary, the numerical methods described
here are based on a finite sequence of operations and lead consequently to a
numerical solution which differs from the exact solution, but agrees with the latter
up to some desired term. We consider therefore:
• a Taylor-series expansion about the current value ti of the independent variable t,
and
• an expansion about ti which is associated with the numerical method used and
differs from the Taylor-series expansion only in the values of its constant
coefficients.
According to the definition given by Engvall [1], a single-step method is said to

be of order p, if p is the largest value for which the expansion associated with the
method agrees with all terms in the Taylor series through the term containing the nth
derivative. This definition is compatible with that which is given below.
According the definition given by Henrici [2], a single-step method is said to be
of order p, if its increment function
xi þ 1 xi
Uðt; x; hÞ ¼
h
approximates the exact relative increment
h 0 h2 hp1 ðp1Þ
Dðt; x; hÞ ¼ f ðt; xÞ þ f ðt; xÞ þ f 00 ðt; xÞ þ þ f ðt; xÞ þ Oðhp Þ
2! 3! p!
with an error of the order of hp. The big-O notation O(hp) is used in this book to
indicate the error term in an approximation to a function. For example, let us
consider the exponential function
xðtÞ ¼ expðtÞ
and its Taylor-series approximation
1 2
x ðt Þ ¼ 1 þ t þ t þ
2
The terms following ½ t2 in the expansion are said to be O(t3), by writing
1 2
expðtÞ ¼ 1 þ t þ t þ O t3
2
to mean that the error committed in the approximation, that is, the absolute value of
the difference between x(t) and x*(t) is smaller than some constant multiplied by t3
as t tends to zero:
jxðtÞ x ðtÞj\a t3 ðas t ! 0Þ
where a is some nonzero constant (a 6¼ 0). In general, some quantity q is said to be

O(hp), whenever
q
lim ¼ constant\1 ðconstant 6¼ 0Þ
h!0 hp
Following Henrici [2], the largest positive integer p such that
U D ¼ O ð hp Þ
is called the order of a single-step method. In general, the higher the order of a
numerical method, the more accurate it is. The order p of the method may be
considered as a first rough measure of the accuracy of the method, because p states
that the local truncation error ei+1 defined as follows
ei þ 1 ¼ hfU½ti ; xðti Þ; h D½ti ; xðti Þ; hg
is comparable to hp+1, but not to any higher power of h.
6.2 Fundamental Concepts on the Runge–Kutta Methods
This group of methods is due to the German mathematicians Carl Runge and Martin
Kutta. The fourth-order Runge–Kutta methods are widely used in computer solu-
tion of differential equations, particularly when the calculation of derivatives of
order higher than the first would be a complex task. These methods will be shown
below for the solution of a first-order differential equation, but an equation of higher
order may easily be written as a system of first-order equations, as has been shown
in Sect. 6.1. The present paragraph gives intuitive concepts on the methods of the
Runge–Kutta family.
Let us consider a differential equation x′= f(t, x) with an associated initial condition
x(t0) = x0. The fundamental idea, on which all Runge–Kutta methods are based, is to
express the increment xn+1 − xn as the step size h multiplied by a weighted average of
certain values k1, k2, …, ks computed at certain points located at abscissae, respec-
tively, (tn + c1h), (tn + c2h), …, (tn + csh) between the beginning and the end of the
current interval [tn, tn+1].
In order to illustrate this concept, we consider first the simplest of these methods,
which performs two evaluations (k1 and k2) of the integrand function f(t, x). These
two evaluations (also called stages) occur at, respectively, the beginning (tn) and the
end (tn+1) of each interval [tn, tn+1]. Let us consider the following two-stage Runge–
Kutta method
x n þ 1 ¼ x n þ hð b1 k 1 þ b2 k 2 Þ
where
k1 ¼ f ð t n ; xn Þ
k2 ¼ f ðtn þ c2 h; xn þ ha21 k1 Þ
and b1, b2, c2, and a21 are certain coefficients to be determined. If we choose
1
b1 ¼ b 2 ¼ c2 ¼ 1 a21 ¼ 1
2
then we have
1
xn þ 1 ¼ xn þ hð k 1 þ k 2 Þ
2
k2 ¼ f ðtn þ h; xn þ hk1 Þ ¼ f ½tn þ h; xn þ hf ðtn ; xn Þ
By so doing, the integrand function f(t, x) is evaluated two times: at the beginning
(t = tn) and the end (t = tn+1) of each integration step; in addition, equal weights
(b1 = b2 = ½) are assigned to these two evaluations (k1 and k2).
In the classical Runge–Kutta formulae, four values k1, k2, k3, and k4 are defined
as follows
k1 ¼ f ðtn ; xn Þ
k2 ¼ f ðtn þ c2 h; xn þ ha21 k1 Þ
k3 ¼ f ½tn þ c3 h; xn þ hða31 k1 þ a32 k2 Þ
k4 ¼ f ½tn þ c4 h; xn þ hða41 k1 þ a42 k2 þ a43 k3 Þ
where
c2 ¼ a21
c3 ¼ a31 þ a32
c4 ¼ a41 þ a42 þ a43
A weighted average U of the values k1, k2, k3, and k4 is computed by means of four
weights b1, b2, b3, and b4, as follows
U ¼ b1 k 1 þ b2 k 2 þ b3 k 3 þ b4 k 4
The definition given above shows that k1 is the first estimate of the angular coef-
ficient of the tangent to the curve at t = tn, that is, at the initial point of the interval
[tn, tn + h]; the subsequent values k2, k3, and k4 are further estimates of the angular
coefficient when t is equal to tn plus some fractions of the step size h; finally, xn+1 at
t = tn+1 is computed as xn plus h multiplied by the weighted average U = b1k1 +
b2k2 + b3k3 + b4k4 of the angular coefficients computed previously.
6.2 Fundamental Concepts on the Runge–Kutta Methods 863
According to a geometrical interpretation of the method, we compute the slope

k1 of the tangent at the initial point (tn, xn) of the step. Using this value of the slope,
we go forward by a fraction c2 of the step size and compute there the slope k2. Now,
using k2 as a new value of the slope, we start again at (tn, xn), go forward by a
fraction c3 of the step size, and compute there the slope k3. Using k3 as the slope, we
again start at (tn, xn), but this time we go a full step forward (c4 = 1) to compute the
slope k4. The weighted average U of the values k1, k2, k3, and k4 is multiplied by
h and the product hU provides the increment that must be added to xn to compute
xn+1. The free coefficients in the expressions written above are ten: the four external
weights b1, b2, b3, and b4 plus the six internal weights a21, a31, a32, a41, a42, and
a43. The nodes c2, c3, and c4, that is, the fractions of the step size h, are functions of
the internal weights, as follows
c2 ¼ a21
c3 ¼ a31 þ a32
c4 ¼ a41 þ a42 þ a43
Such coefficients must be chosen in such a way as to make xi + hU the best possible
approximation to x(tn + h), where x(tn + h) is expressed as a Taylor-series expan-
sion up to the fourth order. In other words, there must be agreement through the
fourth order between xn + hU and the Taylor-series expansion of the exact solution
x(tn + h) around tn. As has been shown in Sect. 6.1, this means that the local
truncation error en+1 of this method must be O(h5). This requirement leads to the
following system of eight algebraic equations
b1 þ b2 þ b3 þ b4 ¼ 1 b2 c2 þ b3 c3 þ b4 c4 ¼ 12
b2 c22 þ b3 c23 þ b4 c24 ¼ 13 b3 c2 a32 þ b4 ðc3 a43 þ c2 a42 Þ ¼ 16
b2 c32 þ b3 c33 þ
b24 c4 ¼ 42
3 1
b3 c2 c3 a32 þ b4 c4 ðc3 a43 þ c2 a42 Þ ¼ 18
b3 c2 a32 þ b4 c3 a43 þ c2 a42 ¼ 12
2 1
b4 c2 a32 a43 ¼ 24
1
Here, the unknowns are the ten free coefficients, so that the values of two of them
may be chosen arbitrarily. If we choose c2 = c3, then
c2 ¼ c3 ¼ 12 c4 ¼ 1 b1 ¼ b4 ¼ 16 b2 þ b3 ¼ 23
a42 þ a43 ¼ 1 b3 a32 ¼ 16 a32 a43 ¼ 12
If we further choose b2 = b3, then b2 = b3 = 1=3, and therefore,
c2 ¼ c3 ¼ 12 a42 ¼ 0 b1 ¼ b4 ¼ 16 c4 ¼ 1
a43 ¼ 1 b2 ¼ b3 ¼ 3 a32 ¼ 12
1
Substituting these values in the definitions given above, we have (Runge’s

formulae)

1 1
k2 ¼ f tn þ h; xn þ h k1
2 2

1 1
2 2
k4 ¼ f ðtn þ h; xn þ h k3 Þ
1
U ¼ ðk1 þ 2k2 þ 2k3 þ k4 Þ
6
1
xn þ 1 ¼ xn þ h U ¼ xn þ hðk1 þ 2k2 þ 2k3 þ k4 Þ
6
Alternatively, if we choose c2 = 1=3 and c3 = 1 − c2 = 2=3, then we have
a42 ¼ 1 b1 ¼ b4 ¼ 18 c4 ¼ 1
a43 ¼ 1 b2 ¼ b3 ¼ 38 a32 ¼ 1
and hence (Kutta’s formulae)

1 1
3 3

2 1
k3 ¼ f tn þ h; xn h k1 þ h k2
3 3
k4 ¼ f ðtn þ h; xn þ h k1 h k2 þ h k3 Þ
1
U ¼ ðk1 þ 3k2 þ 3k3 þ k4 Þ
8
1
8
It is to be noted that c4 = 1 is a consequence of the system of algebraic equations,

independently of the parameters chosen arbitrarily.
As an example, due to Henrici [2], of application of the four-stage Runge–Kutta
formulae, let us consider the following differential equation
x0 ðt; xÞ ¼ t x2
6.2 Fundamental Concepts on the Runge–Kutta Methods 865
with the following initial condition
x ð 0Þ ¼ 0
Let us perform two steps of numerical integration, with a step size h = 0.1, using
the following fourth-order Runge formulae

1 1
2 2

1 1
k3 ¼ f tn þ h; xn þ hk2
2 2
k4 ¼ f ðtn þ h; xn þ h k3 Þ
1
U ¼ ðk1 þ 2k2 þ 2k3 þ k4 Þ
6
1
6
First step (from t = 0.0 to t = 0.1)
k1 ¼ 0:0 ð0:0Þ2 ¼ 0:0

k2 ¼ 0:0 þ 0:1=2 ð0:0 þ 0:1 0:0=2Þ2 ¼ 0:05
k3 ¼ 0:0 þ 0:1=2 ð0:0 þ 0:1 0:05=2Þ2 ¼ 0:04999375
k4 ¼ 0:0 þ 0:1 ð0:0 þ 0:1 0:04999375Þ2 ¼ 0:099975006
U ¼ ð0:0 þ 2 0:05 þ 2 0:04999375 þ 0:099975006Þ=6 ¼ 0:049993751
xð0:1Þ ¼ 0:0 þ 0:1 0:049993751 ¼ 0:004999375
Second step (from t = 0.1 to t = 0.2)
k1 ¼ 0:1ð0:004999375Þ2 ¼ 0:099975006
k2 ¼ 0:1 þ 0:1=2ð0:004999375 þ 0:099975006 0:1=2Þ2 ¼ 0:149900037
k3 ¼ 0:1 þ 0:1=2ð0:004999375 þ 0:149900037 0:1=2Þ2 ¼ 0:149843891
k4 ¼ 0:1 þ 0:1ð0:004999375 þ 0:149843891Þ2 ¼ 0:199600649
U ¼ ð0:099975006 þ 2 0:149900037 þ 2 0:149843891 þ 0:199600649Þ=6
¼ 0:149843919
xð0:2Þ ¼ 0:004999375 þ 0:1 0:149843919 ¼ 0:019983767
In case of Runge–Kutta methods having order higher than four, the complexity of
the equations increases considerably with the order p of the method.
6.3 Runge–Kutta Fourth-Order Methods with Local

Truncation Error Control
For the Runge–Kutta-type methods described above, it is difficult to find a pro-

cedure for step-size control, because they give no information on the value of the
truncation error at each step. For this purpose, it is possible to choose the four stages
k1, k2, k3, and k4 of a classical fourth-order Runge–Kutta formula so that the local
truncation error ei+1, relating to the step [ti, ti+1], can be made equal to
C5h5x(V)(n) + O(h6), for some known coefficient C5, where n is contained in the
interval [ti, ti+1]. Likewise, a linear combination of the stages themselves can also be
chosen to equal C5h5x(V)(n) + O(h6), on a certain class of functions, for example,
linear functions. This makes it possible to have not only an approximate solution
but also an estimate of its local truncation error. For example, by using the classical
fourth–order Runge formula of Sect. 6.2
k1 ¼ f ðti ; xi Þ

1 1
k2 ¼ f ti þ h; xi þ h k1
2 2

1 1
2 2
k4 ¼ f ðti þ h; xi þ h k3 Þ
1
U ¼ ðk1 þ 2k2 þ 2k3 þ k4 Þ
6
1
xi þ 1 ¼ xi þ h U ¼ xi þ hðk1 þ 2k2 þ 2k3 þ k4 Þ
6
and by evaluating a fifth additional stage

3 1
k3 ¼ f ti þ h; xi þ hð5k1 þ 7k2 þ 13k3 k4 Þ
4 32
the local truncation error of this formula is found to be approximately
2
i þ 1 hðk1 þ 3k2 þ 3k3 þ 3k4 8k5 Þ þ O h6
3
Consequently, the computation of a further stage (k5) provides the numerical

solution and a rough estimate of its local truncation error.
6.3 Runge–Kutta Fourth-Order Methods with Local Truncation Error Control 867
Following this line of reasoning, Merson proposed to perform an additional

evaluation (k5) of the integrand function f(t, x) at each step. As has been shown
above, this makes it possible to estimate the local truncation error. As a result of this
estimate, the step size relating to the interval [ti, ti+1] may either be kept constant or
changed, as necessary. The formulae proposed by Merson are

1 1
3 3

1 1 1
k3 ¼ f ti þ h; xi þ h k1 þ h k2
3 6 6

1 1 3
k4 ¼ f ti þ h; xi þ h k1 þ h k3
2 8 8

1 3
k5 ¼ f ti þ h; xi þ h k1 h k3 þ 2h k4
2 2
1
U ¼ ðk1 þ 4k4 þ k5 Þ
6
1
xi þ 1 ¼ xi þ h U ¼ xi þ hðk1 þ 2k4 þ k5 Þ
6
The local truncation error in the formulae given above is approximately
h
i þ 1 ð2k1 9k3 þ 8k4 k5 Þ
30
which tends to zero with h. The decision on whether the step size should, or should
not, be changed is taken on the basis of the value of the local truncation error.
Rogers [3] suggests the following “dipstick” procedure:
• if the local truncation error ei+1 lies outside given limits (one-eighth and twice
the pre-assigned value), then the step size is either doubled or halved;
• otherwise, the step size is left unchanged.
A disadvantage of the Runge–Kutta–Merson method is the fact that the estimate
of the local truncation error is asymptotically correct for linear differential equa-
tions. In the nonlinear case, the method may lead to highly inefficient calculations
[3].
Alternatively, a method suggested by Fehlberg [4] is based on a pair of Runge–
Kutta formulae which use the same set of coefficients but provide two approxi-
mations of different orders. For example, the Runge–Kutta–Fehlberg RK4(5)6
formulae of the fourth order are based on two estimates (one with order of accuracy
four and one with order of accuracy five) of the new value xi+1. The two estimates
xi þ 1 ½4 and xi þ 1 ½5 , where the superscript indicates the order of the method, use the
same evaluations (k1, k2, …, k6) of the integrand function at each step. The Runge–
Kutta–Fehlberg formulae [4, p. 13] are given below.
k1 ¼ f ð t i ; xi Þ

1 1
k2 ¼ f ti þ h; xi þ hk1
4 4

3 3 9
k3 ¼ f ti þ h; xi þ hk1 þ hk2
8 32 32

12 1932 7200 7296
k4 ¼ f ti þ h; xi þ hk1 hk2 þ hk3
13 2197 2197 2197

439 3680 845
k5 ¼ f ti þ h; xi þ hk1 8hk2 þ hk3 hk4
216 513 4104

1 8 3544 1859 11
k6 ¼ f ti þ h; xi hk1 þ 2hk2 hk3 þ hk4 hk5
2 27 2565 4104 40

25 1408 2197 1
xi þ 1 ½4 ¼ xi þ h k1 þ k3 þ k4 k5
216 2565 4104 5

16 6656 28561 9 2
xi þ 1 ½5 ¼ xi þ h k1 þ k3 þ k4 k5 þ k6
135 12825 56430 50 55

1 128 2197 1 2
i þ 1 h k1 k3 k4 þ k5 þ k6
360 4275 75240 50 55
The local truncation error (ei+1) results from the absolute value of the difference
of the two estimates and provides the means of determining the next step size (hi+1)
from the present step size (hi) and the given tolerance (e).
By approximating the true solution x(ti+1) to the higher-order estimate, that is, by
setting xðti þ 1 Þ x½5 i þ 1 , we compute the local truncation error from

ei þ 1 ðhi Þ
x½5 i þ 1 x½4 i þ 1
If this value is greater than the tolerance e, the step must be repeated with a
smaller step-size hi+1. Since the local truncation error is proportional to h5i for the
fourth-order formula, then the local truncation error with hi+1 will be

hi þ 1 5
½5 ½4

hi þ 1 5

i þ 1 ðhi þ 1 Þ ¼ i þ 1 ðhi Þ x iþ1 x iþ1

hi hi
By requiring ei+1(hi+1) to be smaller than ei+1(hi) and solving for hi+1, we have
15

hi þ 1 ¼ hi
i þ 1 ðhi Þ
6.3 Runge–Kutta Fourth-Order Methods with Local Truncation Error Control 869
As an example of application, let us consider the following initial-value problem

proposed by Collatz [5]:
t
x0 ¼ x 2 x ð 0Þ ¼ 1
x
The exact solution is x(t) = (2t + 1)½. Let us perform the first step of numerical
integration by using an initial step size h = 0.2 and a tolerance e = 5 10−7. The
Runge–Kutta–Fehlberg formulae given above may be applied as follows:
t ¼ t0 ¼ 0
x ¼ x0 ¼ 1
k1 ¼ x 2t=x ¼ 1 2 0=1 ¼ 1
t ¼ t0 þ h=4 ¼ 0 þ 0:2=4 ¼ 0:05

x ¼ x0 þ hk1 =4 ¼ 1 þ 0:2 1=4 ¼ 1:05
k2 ¼ x 2t=x ¼ 1:05 2 0:05=1:05 ¼ 0:954761904
t ¼ t0 þ 3h=8 ¼ 0 þ 3 0:2=8 ¼ 0:075

x ¼ x0 þ 3hðk1 þ 3k2 Þ=32 ¼ 1 þ 3 0:2 ð1 þ 3 0:954761904Þ=32
¼ 1:072455357
k3 ¼ x 2t=x ¼ 1:072455357 2 0:075=1:072455357 ¼ 0:932589395
t ¼ t0 þ 12h=13 ¼ 0 þ 12 0:2=13 ¼ 0:184615384

x ¼ x0 þ hð1932k1 7200k2 þ 7296k3 Þ=2197 ¼ 1 þ 0:2 ð1932 1 7200
0:954761904 þ 7296 0:932589395Þ=2197 ¼ 1:169493538
k4 ¼ x 2t=x ¼ 1:169493538 2 0:184615384=1:169493538 ¼ 0:853775019
t ¼ t0 þ h ¼ 0 þ 0:2 ¼ 0:2
x ¼ x0 þ hð439k1 =2168k2 þ 3680k3 =513845k4 =4104Þ ¼ 1 þ 0:2 ð439 1=216
8 0:954761904 þ 3680 0:932589395=513845 0:853775019=4104Þ
¼ 1:181688551
k5 ¼ x2t=x ¼ 1:1816885512 0:2=1:181688551 ¼ 0:843189883
t ¼ t0 þ h=2 ¼ 0 þ 0:2=2 ¼ 0:1

x ¼ x0 hð8k1 =272k2 þ 3544k3 =25651859k4 =4104 þ 11k5 =40Þ
¼ 10:2 ð8 1=272 0:954761904 þ 3544 0:932589395=2565
1859 0:853775019=4104 þ 11 0:843189883=40Þ ¼ 1:095910061
k6 ¼ x2t=x ¼ 1:0959100612 0:1=1:095910061 ¼ 0:913413333
x½4 ðt0 þ hÞ ¼ x0 þ hð25k1 =216 þ 1408k3 =2565 þ 2197k4 =4104 k5 =5Þ
¼ 1 þ 0:2 ð25 1=216 þ 1408 0:932589395=2565 þ 2197
0:853775019=41040:843189883=5Þ ¼ 1:183215928
x½5 ðt0 þ hÞ ¼ x0 þ hð16k1 =135 þ 6656k3 =12825 þ 28561k4 =564309k5 =50 þ 2k6 =55Þ
¼ 1 þ 0:2 ð16 1=135 þ 6656 0:932589395=12825 þ 28561
0:853775019=564309 0:843189883=50 þ 2 0:913413333=55Þ
¼ 1:183216593
The solution in the same point, computed analytically, is

1
xðt0 þ hÞ ¼ ð2 0:2 þ 1Þ2 ¼ 1:183215957
We compute the local truncation error from

ei þ 1 ðhÞ
x½5 ðt0 þ hÞ x½4 ðt0 þ hÞ
¼ j1:183216593 1:183215928j
¼ 6:65 107
Since this value is greater than the fixed tolerance (6.65 10−7 > 5.0 10−7),
then the result is rejected and the first step must be repeated using a step size h* less
than
15 15
5:0 107
h ¼ 0:2 ¼ 0:1889
i þ 1 ð hÞ 6:65 107
In practice, using a security factor equal to 0.9, the first step may be repeated
with a step size
h ¼ 0:9 0:1889 ¼ 0:17
The formulae shown above hold for a scalar differential equation of the first
order. In the case of a system of n differential equations of the first order, the scalars
x, k1, …, k6, and ei+1 are replaced by corresponding column vectors x, k1, …, k6 and
ei+1; likewise, the scalar-valued function f(t, x) is replaced by a vector-valued
function f(t, x). In this case, some norm of the local truncation error vector ei+1 is
used for step-size control.
6.4 Runge–Kutta Methods with Order Higher Than Four 871
6.4 Runge–Kutta Methods with Order Higher Than Four
These methods provide a step-size control procedure based on the concept of

estimating the leading term of the local truncation error at each step by means of
two approximations (one of order p and the other of order p + 1) to the solution
xi+1. The next step size is selected by means of an expression like that shown in
Sect. 6.3, generalised by using an exponent equal to 1/(p + 1):
p þ1 1

hi þ 1 ¼ hi
i þ 1 ð hi Þ
For example, in the Runge–Kutta–Fehlberg method of the fifth order [6], xi+1 is
computed successively by a fifth-order formula
X
5
xi þ 1 ¼ xi þ h bj fj þ O h6
j¼0
followed by a sixth-order formula used for step-size control
X
7
x i þ 1 ¼ xi þ h bj f j þ O h7
j¼0
where {bj} and {b*j} are weighting coefficients (corresponding to the weights b1, b2, b3,
and b4 in the Runge–Kutta fourth-order formula U = b1k1 + b2k2 + b3k3 + b4k4) and

P
j1
fj ¼ f ti þ cj h; xi þ h aj k f k ðj ¼ 1; 2; . . .; 7Þ
k¼0
The coefficients {bj} and {b*j} are the same in the expression of xi+1 and x*i+1, so
that both of them can be computed by evaluating eight times the function fj. The
coefficients are given in Table 6.1 (from [6]), where the notation RK5(6)8 means a
Runge–Kutta–Fehlberg method of order 5 with embedded method of order 6 used
for step-size control and 8 evaluations of the integrand function.
The difference between xi+1 and x*i+1 yields the leading term of the local
truncation error:
X
5
5
i þ 1 ¼ h ðbj b j Þfj hðb 6 f6 þ b 7 f7 Þ ¼ hðf0 þ f5 f6 f7 Þ
j¼0
66
As indicated above, the principle on which the Runge–Kutta–Fehlberg methods are

based is to compute two Runge–Kutta estimates for the new value xi, having
different order of error (e.g., fifth and sixth order in the formulae given above). At
each step, the leading term of the local truncation error is evaluated, and the next
step size can be computed according to the value of the error.
Table 6.1 Coefficients for RK5(6)8

c0 = 0
c1 = 1/6
c2 = 4/15
c3 = 2/3
c4 = 4/5
c5 = 1
c6 = 0
c7 = 1
a10 = 1/6
a20 = 4/75
a21 = 16/75
a30 = 5/6
a31 = −8/3
a32 = 5/2
a40 = −8/5
a41 = 144/25
a42 = −4
a43 = 16/25
a50 = 361/320
a51 = −18/5
a52 = 407/128
a53 = −11/80
a54 = 55/128
a60 = −11/640
a61 = 0
a62 = 11/256
a63 = −11/160
a64 = 11/256
a65 = 0
a70 = 93/640
a71 = −18/5
a72 = 803/256
a73 = −11/160
a74 = 99/256
a75 = 0
a76 = 1
b0 = 31/384
b1 = 0
b2 = 1125/2816
b3 = 9/32
b4 = 125/768
b5 = 5/66
b*0 = 7/1408
b*1 = 0
b*2 = 1125/2816
b*3 = 9/32
b*4 = 125/768
b*5 = 0
b*6 = 5/66
b*7 = 5/66
In the Runge–Kutta–Fehlberg method of the sixth order [6], xi+1 is computed

successively by a sixth-order formula
X
7
j¼0
followed by a seventh-order formula used for step-size control
X
9
j¼0
where {bj} and {b*j} are weighting coefficients, and

P
j1
fj ¼ f ti þ cj h; xi þ h ajk fk ðj ¼ 1; 2; . . .; 9Þ
k¼0
The coefficients {bj} and {b*j} are computed by evaluating ten times the function
fj. The coefficients are given in Table 6.2 (from [6]).
truncation error:
X
7
11
i þ 1 ¼ h ðbj b j Þfj hðb 8 f8 þ b 9 f9 Þ ¼ hð f 0 þ f 7 f 8 f 9 Þ
j¼0
270

c0 = 0
c1 = 2/33
c2 = 4/33
c3 = 2/11
c4 = 1/2
c5 = 2/3
c6 = 6/7
c7 = 1
c8 = 0
c9 = 1
a10 = 2/33
a20 = 0
a21 = 4/33
a30 = 1/22
a31 = 0
a32 = 3/22
a40 = 43/64
a41 = 0
a42 = −165/64
a43 = 77/32
a50 = −2383/486
(continued)
Table 6.2 (continued)

a51 = 0
a52 = 1067/54
a53 = −26312/1701
a54 = 2176/1701
a60 = 10077/4802
a61 = 0
a62 = −5643/686
a63 = 116259/16807
a64 = −6240/16807
a65 = 1053/2401
a70 = −733/176
a71 = 0
a72 = 141/8
a73 = −335763/23296
a74 = 216/77
a75 = −4617/2816
a76 = 7203/9152
a80 = 15/352
a81 = 0
a82 = 0
a83 = −5445/46592
a84 = 18/77
a85 = −1215/5632
a86 = 1029/18304
a87 = 0
a90 = −1833/352
a91 = 0
a92 = 141/8
a93 = −51237/3584
a94 = 18/7
a95 = −729/512
a96 = 1029/1408
a97 = 0
a98 = 1
b0 = 77/1440
b1 = 0
b2 = 0
b3 = 1771561/6289920
b4 = 32/105
b5 = 243/2560
b6 = 16807/74880
b7 = 11/270
b*0 = 11/864
b*1 = 0
b*2 = 0
b*3 = 1771561/6289920
b*4 = 32/105
b*5 = 243/2560
b*6 = 16807/74880
b*7 = 0
b*8 = 11/270
b*9 = 11/270
In the Runge–Kutta–Fehlberg method of the seventh order [6], xi+1 is computed

successively by a seventh-order formula
X
10
j¼0
followed by an eighth-order formula used for step-size control
X
12
j¼0

P
j1
fj ¼ f ti þ cj h; xi þ h aj k fk ðj ¼ 1; 2; . . .; 12Þ
k¼0
The coefficients {bj} and {b*j} are computed by evaluating thirteen times the
function fj. The coefficients are given in Table 6.3 (from [6]).
truncation error, as follows
X9
41h
i þ 1 ¼ h ðbj b j Þfj hðb 10 f10 þ b 11 f11 Þ ¼ ðf0 þ f10 f11 f12 Þ
j¼0
840

c0 = 0
c1 = 2/27
c2 = 1/9
c3 = 1/6
c4 = 5/12
c5 = 1/2
c6 = 5/6
c7 = 1/6
c8 = 2/3
c9 = 1/3
c10 = 1
c11 = 0
c12 = 1
a10 = 2/27
a20 = 1/36
a21 = 1/12
a30 = 1/24
a31 = 0
a32 = 1/8
(continued)

a40 = 5/12
a41 = 0
a42 = −25/16
a43 = 25/16
a50 = 1/20
a51 = 0
a52 = 0
a53 = 1/4
a54 = 1/5
a60 = −25/108
a61 = 0
a62 = 0
a63 = 125/108
a64 = −65/27
a65 = 125/54
a70 = 31/300
a71 = 0
a72 = 0
a73 = 0
a74 = 61/225
a75 = −2/9
a76 = 13/900
a80 = 2
a81 = 0
a82 = 0
a83 = −53/6
a84 = 704/45
a85 = −107/9
a86 = 67/90
a87 = 3
a90 = −91/108
a91 = 0
a92 = 0
a93 = 23/108
a94 = −976/135
a95 = 311/54
a96 = −19/60
a97 = 17/6
a98 = −1/12
a100 = 2383/4100
a101 = 0
a102 = 0
a103 = −341/164
a104 = 4496/1025
a105 = −301/82
a106 = 2133/4100
a107 = 45/82
a108 = 45/164
a109 = 18/41
a110 = 3/205
a111 = 0
(continued)

a112 = 0
a113 = 0
a114 = 0
a115 = −6/41
a116 = −3/205
a117 = −3/41
a118 = 3/41
a119 = 6/41
a1110 = 0
a120 = −1777/4100
a121 = 0
a122 = 0
a123 = −341/164
a124 = 4496/1025
a125 = −289/82
a126 = 2193/4100
a127 = 51/82
a128 = 33/164
a129 = 12/41
a1210 = 0
a1211 = 1
b0 = 41/840
b1 = 0
b2 = 0
b3 = 0
b4 = 0
b5 = 34/105
b6 = 9/35
b7 = 9/35
b8 = 9/280
b9 = 9/280
b10 = 41/840
b*0 = 0
b*1 = 0
b*2 = 0
b*3 = 0
b*4 = 0
b*5 = 34/105
b*6 = 9/35
b*7 = 9/35
b*8 = 9/280
b*9 = 9/280
b*10 = 0
b*11 = 41/840
b*12 = 41/840
In the Runge–Kutta–Fehlberg method of the eighth order [6], xi+1 is computed

successively by an eighth-order formula
X
14
j¼0
followed by a ninth-order formula used for step-size control
X
16
x i þ 1 ¼ xi þ h bj fj þ O h10
j¼0

P
j1
k¼0
The coefficients {bj} and {b*j} are computed by evaluating seventeen times the
function fj. The coefficients are given in Table 6.4 (from [6]), where the coefficients
ajk and bj not listed are zero. In addition, b*8 = b8, b*9 = b9, …, b*13 = b13,
b*14 = 0, b*15 = b*16 = b14.
The leading term of the local truncation error is
ei þ 1 ¼ b14 ðf0 þ f14 f15 f16 Þh

c1 = 0.44368940376498183109599404281370
c2 = 0.66553410564747274664399106422055
c3 = 0.99830115847120911996598659633083
c4 = 0.31550000000000000000000000000000
c5 = 0.50544100948169068626516126737384
c6 = 0.17142857142857142857142857142857
c7 = 0.82857142857142857142857142857143
c8 = 0.66543966121011562534953769255586
c9 = 0.24878317968062652069722274560771
c10 = 0.10900000000000000000000000000000
c11 = 0.89100000000000000000000000000000
c12 = 0.39950000000000000000000000000000
c13 = 0.60050000000000000000000000000000
c14 = 1
c15 = 0
c16 = 1
a10 = 0.44368940376498183109599404281370
a20 = 0.16638352641186818666099776605514
a21 = 0.49915057923560455998299329816541
a30 = 0.24957528961780227999149664908271
a32 = 0.74872586885340683997448994724812
(continued)

a40 = 0.20661891163400602426556710393185
a42 = 0.17707880377986347040380997288319
a43 = −0.68197715413869494669377076815048 10−1
a53 = 0.40215962642367995421990563690087 10−2
a54 = 0.39214118169078980444392330174325
a60 = 0.98899281409164665304844765434355 10−1
a63 = 0.35138370227963966951204487356703 10−2
a64 = 0.12476099983160016621520625872489
a65 = −0.55745546834989799643742901466348 10−1
a70 = −0.36806865286242203724153101080691
a74 = −0.22273897469476007645024020944166 10+1
a75 = 0.13742908256702910729565691245744 10+1
a76 = 0.20497390027111603002159354092206 10+1
a80 = 0.45467962641347150077351950603349 10−1
a85 = 0.32542131701589147114677469648853
a86 = 0.28476660138527908888182420573687
a87 = 0.97837801675979152435868397271099 10−2
a90 = 0.60842071062622057051094145205182 10−1
a95 = −0.21184565744037007526325275251206 10−1
a96 = 0.19596557266170831957464490662983
a97 = −0.42742640364817603675144835342899 10−2
a98 = 0.17434365736814911965323452558189 10−1
a10,0 = 0.54029783296931917365785724111182 10−1
a10,6 = 0.11029825597828926530283127648228
a10,7 = −0.12565008520072556414147763782250 10−2
a10,8 = 0.36790043477581460136384043566339 10−2
a10,9 = −0.57780542770972073040840628571866 10−1
a11,0 = 0.12732477068667114646645181799160
a11,7 = 0.11448805006396105323658875721817
a11,8 = 0.28773020709697992776202201849198
a11,9 = 0.50945379459611363153735885079465
a11,10 = −0.14799682244372575900242144449640
a12,0 = −0.36526793876616740535848544394333 10−2
a12,5 = 0.81629896012318919777819421247030 10−1
a12,6 = −0.38607735635693506490517694343215
a12,7 = 0.30862242924605106450474166025206 10−1
a12,8 = −0.58077254528320602815829374733518 10−1
a12,9 = 0.33598659328884971493143451362322
a12,10 = 0.41066880401949958613549622786417
a12,11 = −0.11840245972355985520633156154536 10−1
a13,0 = −0.12375357921245143254979096135669 10+1
a13,5 = −0.24430768551354785358734861366763 10+2
a13,6 = 0.54779568932778656050436528991173
a13,7 = −0.44413863533413246374956896569346 10+1
a13,8 = 0.10013104813713266094792617851022 10+2
a13,9 = −0.14995773102051758447170985073142 10+2
a13,10 = 0.58946948523217013620824539651427 10+1
a13,11 = 0.17380377503428984877616857440542 10+1
(continued)

a13,12 = 0.27512330693166730263758622860276 10+2
a14,0 = −0.35260859388334522700502958875588
a14,5 = −0.18396103144848270375044198988231
a14,6 = −0.65570189449741645138006879985251
a14,7 = −0.39086144880439863435025520241310
a14,8 = 0.26794646712850022936584423271209
a14,9 = −0.10383022991382490865769858507423 10+1
a14,10 = 0.16672327324258671664727346168501 10+1
a14,11 = 0.49551925855315977067732967071441
a14,12 = 0.11394001132397063228586738141784 10+1
a14,13 = 0.51336696424658613688199097191534 10−1
a15,0 = 0.10464847340614810391873002406755 10−2
a15,8 = −0.67163886844990282237778446178020 10−2
a15,9 = 0.81828762189425021265330065248999 10−2
a15,10 = −0.42640342864483347277142138087561 10−2
a15,11 = 0.28009029474168936545976331153703 10−3
a15,12 = −0.87835333876238676639057813145633 10−2
a15,13 = 0.10254505110825558084217769664009 10−1
a16,0 = −0.13536550786174067080442168889966 10+1
a16,5 = −0.18396103144848270375044198988231
a16,6 = −0.65570189449741645138006879985251
a16,7 = −0.39086144880439863435025520241310
a16,8 = 0.27466285581299925758962207732989
a16,9 = −0.10464851753571915887035188572676 10+1
a16,10 = 0.16714967667123155012004488306588 10+1
a16,11 = 0.49523916825841808131186990740287
a16,12 = 0.11481836466273301905225795954930 10+1
a16,13 = 0.41082191313833055603981327527525 10−1
a16,15 = 1
b0 = 0.32256083500216249913612900960247 10−1
b8 = 0.25983725283715403018887023171963
b9 = 0.92847805996577027788063714302190 10−1
b10 = 0.16452339514764372891647731842800
b11 = 0.17665951637860074367084298397547
b12 = 0.23920102320352759374108933320941
b13 = 0.39484274604202853746752118829325 10−2
b14 = 0.30726495475860640406368305522124 10−1
As Verner [7] has shown, the Runge–Kutta–Fehlberg formula pairs 5(6)8, 6(7)10, 7
(8)13, and 8(9)16 have error estimators which give erroneously zero values for
quadrature problems, that is, for problems in which the integrand function f(t, x)
depends only on the independent variable t:
x0 ðtÞ ¼ f ðtÞ with xðt0 Þ ¼ x0

This is because, in these cases, the method of order p has exactly the same trun-
cation error as the method of order p + 1. Consequently, the solution found at each
step will be accepted (and possibly the step-size increased), without regard to the
magnitude of the local truncation error [7]. For the purpose of removing this defect,
Verner has derived other Runge–Kutta pairs than those previously found by
Fehlberg. The results obtained by Verner are shown below. The notation, indicating
the two consecutive orders and the number of stages for each formula pair, is the
same as that adopted above; now, RKV stands for Runge–Kutta–Verner.
In the Runge–Kutta–Verner method of the fifth order [7], xi+1 is computed
successively by a fifth-order formula
X
6
j¼1
followed by a sixth-order formula
X
8
j¼1

P
j1
k¼1
The leading term of the local truncation error is expressed by

33f1 132f3 891f4 33f5 73f6 891f7 2f8
i þ 1 ¼h þ þ þ
640 325 2240 320 700 8320 35
Here, as is the case with Fehlberg’s formulae, we can choose one of two options:
(a) using the fifth-order formula to compute the solution and the sixth-order formula
only for step-size control; or (b) doing the reverse. Fehlberg [6] chooses option (a),
as has been shown above. Butcher [8] advises to choose option (b), that is, the
sixth-order formula for the solution and the fifth-order formula only for step-size
control, because the sixth-order formula has a larger stability region than that of the
fifth-order formula.
The coefficients are given in Table 6.5 (from [7]).
Table 6.5 Coefficients for RKV5(6)8

c1 = 0
c2 = 1/18
c3 = 1/6
c4 = 2/9
c5 = 2/3
c6 = 1
c7 = 8/9
c8 = 1
a21 = 1/18
a31 = −1/12
a32 = 1/4
a41 = −2/81
a42 = 4/27
a43 = 8/81
a51 = 40/33
a52 = −4/11
a53 = −56/11
a54 = 54/11
a61 = −369/73
a62 = 72/73
a63 = 5380/219
a64 = −12285/584
a65 = 2695/1752
a71 = −8716/891
a72 = 656/297
a73 = 39520/891
a74 = −416/11
a75 = 52/27
a76 = 0
a81 = 3015/256
a82 = −9/4
a83 = −4219/78
a84 = 5985/128
a85 = −539/384
a86 = 0
a87 = 693/3328
b1 = 3/80
b2 = 0
b3 = 4/25
b4 = 243/1120
b5 = 77/160
b6 = 73/700
b*1 = 57/640
b*2 = 0
b*3 = −16/65
b*4 = 1377/2240
b*5 = 121/320
b*6 = 0
b*7 = 891/8320
b*8 = 2/35
In the Runge–Kutta–Verner method of the sixth order [7], xi+1 is computed

successively by a sixth-order formula
X
8
j¼1
followed by a seventh-order formula
X
10
j¼1

P
j1
k¼1

17f1 272f4 272f5 24137569f6 34f7 7f8
i þ 1 ¼ h þ þ
2688 4935 273 57482880 105 90

4131f9 157f10
þ
3920 1260

c1 = 0
c2 = 1/12
c3 = 1/6
c4 = 1/4
c5 = 3/4
c6 = 16/17
c7 = 1/2
c8 = 1
c9 = 2/3
c10 = 1
a21 = 1/12
a31 = 0
a32 = 1/6
a41 = 1/16
a42 = 0
a43 = 3/16
a51 = 21/16
a52 = 0
a53 = −81/16
a54 = 9/2
a61 = 1344688/250563
a62 = 0
(continued)

a63 = −1709184/83521
a64 = 1365632/83521
a65 = −78208/250563
a71 = −559/384
a72 = 0
a73 = 6
a74 = −204/47
a75 = 14/39
a76 = −4913/78208
a81 = −625/224
a82 = 0
a83 = 12
a84 = −456/47
a85 = 48/91
a86 = 14739/136864
a87 = 6/7
a91 = −12253/99144
a92 = 0
a93 = 16/27
a94 = 16/459
a95 = 29072/161109
a96 = −2023/75816
a97 = 112/12393
a98 = 0
a10,1 = 30517/2512
a10,2 = 0
a10,3 = −7296/157
a10,4 = 268728/7379
a10,5 = 2472/2041
a10,6 = −3522621/10743824
a10,7 = 132/157
a10,8 = 0
a10,9 = −12393/4396
b1 = 7/90
b2 = 0
b3 = 0
b4 = 16/45
b5 = 16/45
b6 = 0
b7 = 2/15
b8 = 7/90
b*1 = 2881/40320
b*2 = 0
b*3 = 0
b*4 = 1216/2961
b*5 = −2624/4095
b*6 = 24137569/57482880
b*7 = −4/21
b*8 = 0
b*9 = 4131/3920
b*10 = −157/1260
In the Runge–Kutta–Verner method of the seventh order [7], xi+1 is computed

successively by a seventh-order formula
X
11
j¼1
followed by an eighth-order formula
X
13
j¼1

P
j1
k¼0

f1 16f6 2401f7 2401f8 243f9 2401f10 19f11
i þ 1 ¼ h þ þ
480 375 528000 132000 14080 19200 450

243f12 31f13
þ þ
1760 720

c1 = 0
c2 = 1/4
c3 = 1/12
c4 = 1/8
c5 = 2/5
c6 = 1/2
c7 = 6/7
c8 = 1/7
c9 = 2/3
c10 = 2/7
c11 = 1
c12 = 1/3
c13 = 1
a21 = 1/4
a31 = 5/72
a32 = 1/72
a41 = 1/32
a42 = 0
a43 = 3/32
(continued)

a51 = 106/125
a52 = 0
a53 = −408/125
a54 = 352/125
a61 = 1/48
a62 = 0
a63 = 0
a64 = 8/33
a65 = 125/528
a71 = −1263/2401
a72 = 0
a73 = 0
a74 = 39936/26411
a75 = −64125/26411
a76 = 5520/2401
a81 = 37/392
a82 = 0
a83 = 0
a84 = 0
a85 = 1625/9408
a86 = −2/15
a87 = 61/6720
a91 = 17176/25515
a92 = 0
a93 = 0
a94 = −47104/25515
a95 = 1325/504
a96 = −41792/25515
a97 = 20237/145800
a98 = 4312/6075
a10,1 = −23834/180075
a10,2 = 0
a10,3 = 0
a10,4 = −77824/1980825
a10,5 = −636635/633864
a10,6 = 254048/300125
a10,7 = −183/7000
a10,8 = 8/11
a10,9 = −324/3773
a11,1 = 127/7600
a11,2 = 0
a11,3 = 0
a11,4 = −20032/5225
a11,5 = 456485/80256
a11,6 = −42599/7125
a11,7 = 339227/912000
a11,8 = −1029/4180
a11,9 = 1701/1408
a11,10 = 5145/2432
a12,1 = −27061/204120
a12,2 = 0
(continued)

a12,3 = 0
a12,4 = 40448/280665
a12,5 = −1353775/1197504
a12,6 = 17662/25515
a12,7 = −71687/1166400
a12,8 = 98/225
a12,9 = 1/16
a12,10 = 3773/11664
a12,11 = 0
a13,1 = 11203/8680
a13,2 = 0
a13,3 = 0
a13,4 = −38144/11935
a13,5 = 2354425/458304
a13,6 = −84046/16275
a13,7 = 673309/1636800
a13,8 = 4704/8525
a13,9 = 9477/10912
a13,10 = −1029/992
a13,11 = 0
a13,12 = 729/341
b1 = 13/288
b2 = 0
b3 = 0
b4 = 0
b5 = 0
b6 = 32/125
b7 = 31213/144000
b8 = 2401/12,375
b9 = 1701/14080
b10 = 2401/19200
b11 = 19/450
b*1 = 31/720
b*2 = 0
b*3 = 0
b*4 = 0
b*5 = 0
b*6 = 16/75
b*7 = 16807/79200
b*8 = 16807/79200
b*9 = 243/1760
b*10 = 0
b*11 = 0
b*12 = 243/1760
b*13 = 31/720
In the Runge–Kutta–Verner method of the eighth order [7], xi+1 is computed

successively by an eighth-order formula
X
14
j¼1
followed by a ninth-order formula
X
16
x i þ 1 ¼ xi þ h bj fj þ O h10
j¼1

P
j1
k¼1

7f1 63f8 14f9 21f10 1024f11 21f12 3f13
i þ 1 ¼h þ þ
400 200 25 20 975 36400 25

9f14 9f15 233f16
þ þ
280 25 4200

c1 = 0
c2 = 1/12
c3 = 1/9
c4 = 1/6
c5 = [2 + 2(6)1/2]/15
c6 = [6 + (6)1/2]/15
c7 = [6 − (6)1/2]/15
c8 = 2/3
c9 = 1/2
c10 = 1/3
c11 = 1/4
c12 = 4/3
c13 = 5/6
c14 = 1
c15 = 1/6
c16 = 1
a21 = 1/12
a31 = 1/27
a32 = 2/27
(continued)

a41 = 1/24
a42 = 0
a43 = 1/8
a51 = [4 + 94(6)1/2]/375
a52 = 0
a53 = [−94 − 84(6)1/2]/125
a54 = [328 + 208(6)1/2]/375
a61 = [9 − (6)1/2]/150
a62 = 0
a63 = 0
a64 = [312 + 32(6)1/2]/1425
a65 = [69 + 29(6)1/2]/570
a71 = [927 − 347(6)1/2]/1250
a72 = 0
a73 = 0
a74 = [−16248 + 179(6)1/2]/9375
a75 = [−489 + 208(6)1/2]/3750
a76 = [14268 − 5798(6)1/2]/9375
a81 = 2/27
a82 = 0
a83 = 0
a84 = 0
a85 = 0
a86 = [16 − (6)1/2]/54
a87 = [16 + (6)1/2]/54
a91 = 19/256
a92 = 0
a93 = 0
a94 = 0
a95 = 0
a96 = [118 − 23(6)1/2]/512
a97 = [118 + 23(6)1/2]/512
a98 = −9/256
a10,1 = 11/144
a10,2 = 0
a10,3 = 0
a10,4 = 0
a10,5 = 0
a10,6 = [266 − (6)1/2]/864
a10,7 = [266 + (6)1/2]/864
a10,8 = −1/16
a10,9 = −8/27
a11,1 = [5034 − 271(6)1/2]/61440
a11,2 = 0
a11,3 = 0
a11,4 = 0
a11,5 = 0
a11,6 = 0
a11,7 = [7859 − 1626(6)1/2]/10240
a11,8 = [−2232 + 813(6)1/2]/20480
a11,9 = [−594 + 271(6)1/2]/960
(continued)

a11,10 = [657 − 813(6)1/2]/5120
a12,1 = [5996 − 3794(6)1/2]/405
a12,2 = 0
a12,3 = 0
a12,4 = 0
a12,5 = 0
a12,6 = [−4342 − 338(6)1/2]/9
a12,7 = [154922 − 40458(6)1/2]/135
a12,8 = [−4176 + 3794(6)1/2]/45
a12,9 = [−340864 + 242816(6)1/2]/405
a12,10 = [26304 − 15176(6)1/2]/45
a12,11 = −26624/81
a13,1 = [3793 + 2168(6)1/2]/103680
a13,2 = 0
a13,3 = 0
a13,4 = 0
a13,5 = 0
a13,6 = [4042 + 2263(6)1/2]/13824
a13,7 = [−231278 + 40717(6)1/2]/69120
a13,8 = [7947 − 2168(6)1/2]/11520
a13,9 = [1048 − 542(6)1/2]/405
a13,10 = [−1383 + 542(6)1/2]/720
a13,11 = 2624/1053
a13,12 = 3/1664
a14,1 = −137/1296
a14,2 = 0
a14,3 = 0
a14,4 = 0
a14,5 = 0
a14,6 = [5642 − 337(6)1/2]/864
a14,7 = [5642 + 337(6)1/2]/864
a14,8 = −299/48
a14,9 = 184/81
a14,10 = −44/9
a14,11 = −5120/1053
a14,12 = −11/468
a14,13 = 16/9
a15,1 = [33617 − 2168(6)1/2]/518400
a15,2 = 0
a15,3 = 0
a15,4 = 0
a15,5 = 0
a15,6 = [−3846 + 31(6)1/2]/13824
a15,7 = [155338 − 52807(6)1/2]/345600
a15,8 = [−12537 + 2168(6)1/2]/57600
a15,9 = [92 + 542(6)1/2]/2025
a15,10 = [−1797 − 542(6)1/2]/3600
a15,11 = 320/567
a15,12 = −1/1920
a15,13 = 4/105
(continued)

a15,14 = 0
a16,1 = [−36487 − 30352(6)1/2]/279600
a16,2 = 0
a16,3 = 0
a16,4 = 0
a16,5 = 0
a16,6 = [−29666 − 4499(6)1/2]/7456
a16,7 = [2779182 − 615973(6)1/2]/186400
a16,8 = [−94329 + 91056(6)1/2]/93200
a16,9 = [−232192 + 121408(6)1/2]/17475
a16,10 = [101226 − 22764(6)1/2]/5825
a16,11 = −169984/9087
a16,12 = −87/30290
a16,13 = 492/1165
a16,14 = 0
a16,15 = 1260/233
b1 = 103/1680
b2 = 0
b3 = 0
b4 = 0
b5 = 0
b6 = 0
b7 = 0
b8 = −27/140
b9 = 76/105
b10 = −201/280
b11 = 1024/1365
b12 = 3/7280
b13 = 12/35
b14 = −9/280
b*1 = 23/525
b*2 = 0
b*3 = 0
b*4 = 0
b*5 = 0
b*6 = 0
b*7 = 0
b*8 = 171/1400
b*9 = 86/525
b*10 = 93/280
b*11 = −2048/6825
b*12 = −3/18,200
b*13 = 39/175
b*14 = 0
b*15 = 9/25
b*16 = 233/4200
6.5 Runge–Kutta–Nyström Methods
The methods described here have been studied by E.J. Nyström in 1925 for special
equations of the second order, where the integrand function does not depend
explicitly on the first derivative of the dependent variable. Such is the case with
Cowell’s and Enke’s formulation of the equations of motion in celestial mechanics,
for which these methods are particularly suited. However, if air drag cannot be
neglected in the force model considered, then first-derivative terms are present in
the integrand function, which fact precludes the use of Runge–Kutta–Nyström
formulae [9].
The advantage of the Nyström methods over those shown above is the possi-
bility of reaching a higher order of agreement with the Taylor-series expansion of
the solution for a given number of evaluations of the derivatives than can be
expected in the general case. Let us consider the special second-order differential
equations written in vector form
x0 ¼ y
y0 ¼ f ðt; xÞ
where x and y are, respectively, the position and velocity vectors, and the prime sign
denotes first derivatives with respect to time. The integrand function f(t, x) does not
depend explicitly on y. The step size h denotes the time interval starting from the
initial time t = t0.
The initial conditions are imposed by specifying the position and velocity
vectors at the time t = t0, that is, x(t0) = x0 and y(t0) = y0.
The general Runge–Kutta–Nyström method of order p (the order of a single-step
method has been defined in Sect. 6.1) with s stages (i.e., evaluations of the inte-
grand function at each step) has the following form:
X
s
x i þ 1 ¼ xi þ h y i þ h2 b^ k kk
k¼1
X
s
yi þ 1 ¼ yi þ h bk kk
k¼1
6.5 Runge–Kutta–Nyström Methods 893
where
kk ¼ f ðti ; xi Þ
!
X
s1
ðj ¼ 2; 3; . . .; sÞ
kj ¼ f ti þ hcj ; xi þ hcj yi þ h 2
ajk kk
k¼1
where the order p of the method and the values of the parameters b^k, bk, cj, and ajk
depend on the desired accuracy. This means that they must be chosen so that there
should be an agreement up to the pth order between the expressions given above and
the Taylor-series expansions of x and y for the smallest possible value of the
number of evaluations of the integrand function f(t, x).
Error control is performed by means of two further sets of parameters {b^*k}
and {b*k}, as the next paragraph will show.
Nyström (see [10]) found the following formulae of the fifth order:
First set (c1 = 1=5, c2 = 2=3, and c3 = 1)
1 2
xi þ 1 ¼ xi þ h y i þ h ð14k1 þ 100k2 þ 54k3 Þ þ O h6
336
1
yi þ 1 ¼ yi þ hð14k1 þ 125k2 þ 162k3 þ 35k4 Þ þ O h6
336
where

1 1 1 2
k2 ¼ f ti þ h; xi þ h yi þ h k1
5 5 50

2 2 1
k3 ¼ f ti þ h; xi þ h yi h2 ðk1 7k2 Þ
3 3 27

1 2
k4 ¼ f ti þ h; xi þ h yi þ h ð21k1 4k2 þ 18k3 Þ
70
Second set (c1 = 2=5, c2 = 2=3, and c3 = 4=5)
1 2
xi þ 1 ¼ xi þ h y i þ h ð23k1 þ 75k2 27k3 þ 25k4 Þ þ O h6
192
1
yi þ 1 ¼ yi þ hð23k1 þ 125k2 81k3 þ 125k4 Þ þ O h6
192
where

2 2 2 2
5 5 25

2 2 2
k3 ¼ f ti þ h; xi þ h yi þ h2 k1
3 3 9

4 4 4 2
k4 ¼ f ti þ h; xi þ h yi þ h ð k1 þ k2 Þ
5 5 25
In 1955, Julius Albrecht (see [10]) found the following formulae of the sixth order,
which need five evaluations of the integrand function f(t, x):
1 2
xi þ 1 ¼ xi þ h y i þ h ð7k1 þ 24k2 þ 6k3 þ 8k4 Þ þ O h7
90
1
yi þ 1 ¼ yi þ hð7k1 þ 32k2 þ 12k3 þ 32k4 þ 7k5 Þ þ O h7
90
where

1 1 1 2
4 4 32

1 1 1 2
k3 ¼ f ti þ h; xi þ h yi h ðk1 4k2 Þ
2 2 24

3 3 1 2
k4 ¼ f ti þ h; xi þ h yi þ h ð3k1 þ 4k2 þ 2k3 Þ
4 4 32

1 2
k5 ¼ f ti þ h; xi þ h yi þ h ð6k2 k3 þ 2k4 Þ
14
As an example of application of Albrecht’s formulae, let us consider the following

system of two second-order differential equations, suggested by Bettis [11],
expressed in scalar form
2x2
x00 1 ¼ f1 ðt; x1 ; x2 Þ ¼ 4t2 x1 12
x21 þ x22
2x1
x00 2 ¼ f2 ðt; x1 ; x2 Þ ¼ 4t2 x2 þ 12
x21 þ x22
with the following initial conditions

p12
t0 ¼ ¼ 1:253314137
2
x10 ¼ 0:0
x20 ¼ 1:0
x0 10 y10 ¼ ð2pÞ2 ¼ 2:506628275

1
x0 20 y20 ¼ 0:0
The exact solution is
x1 ðtÞ ¼ cosðt2 Þ
x2 ðtÞ ¼ sinðt2 Þ
In order to apply the method, the first step of numerical integration with a step size
h = 0.01 may be performed as shown below.
t ¼ t0 ¼ 1:253314137
x1 ¼ x10 ¼ 0:0
x2 ¼ x20 ¼ 1:0
1=2
k11 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2533141372 0 2 1= 02 þ 12 ¼ 2:0
1=2
k12 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2533141372 1 þ 2 0= 02 þ 12 ¼ 6:283185304
t ¼ t0 þ h=4 ¼ 1:253314137 þ 0:01=4 ¼ 1:255814137

x1 ¼ x10 þ hy10 =4 þ h2 k11 =32 ¼ 0 þ 0:01 ð2:506628275Þ=4
þ 0:012 ð2:0Þ=32 ¼ 0:00627282
x2 ¼ x20 þ hy20 =4 þ h2 k12 =32 ¼ 1 þ 0:01 0=4 þ 0:012 ð6:283185304Þ=32
¼ 0:999980365
k21 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2558141372 ð0:00627282Þ
1=2
2 0:999980365= 0:006272822 þ 0:9999803652 ¼ 1:960389968
k22 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2558141372 0:999980365
1=2
þ 2 ð0:00627282Þ= 0:006272822 þ 0:9999803652 ¼ 6:320698363
t ¼ t0 þ h=2 ¼ 1:253314137 þ 0:01=2 ¼ 1:258314137

x1 ¼ x10 þ hy10 =2 h2 ðk11 4k21 Þ=24 ¼ 0 þ 0:01 ð2:506628275Þ=2
0:012 ð2 þ 4 1:960389968Þ=24 ¼ 0:012557481
x2 ¼ x20 þ hy20 =2 h2 ðk12 4k22 Þ=24 ¼ 1 þ 0:01 0=2
0:012 ð6:283185304 þ 4 6:320698363Þ=24 ¼ 0:999920835
k31 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2583141372 ð0:012557481Þ
1=2
2 0:999920835= 0:0125574812 þ 0:9999208352 ¼ 1:920310529
k32 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2583141372 0:999920835
1=2
þ 2 ð0:012557481Þ= 0:0125574812 þ 0:9999208352 ¼ 6:358031454
t ¼ t0 þ 3h=4 ¼ 1:253314137 þ 3 0:01=4 ¼ 1:260814137

x1 ¼ x10 þ 3hy10 =4 þ h2 ð3k11 þ 4k21 þ 2k31 Þ=32
¼ 0 þ 3 0:01 ð2:506628275Þ=4 þ 0:012
ð2 31:960389968 41:920310529 2Þ=32 ¼ 0:018854968
x2 ¼ x20 þ 3hy20 =4 þ h2 ð3k12 þ 4k22 þ 2k32 Þ=32 ¼ 1 þ 3 0:01 0=4 þ 0:012
ð6:283185304 36:320698363 4 6:358031454 2Þ=32
¼ 0:999822348
k41 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2608141372 ð0:018854968Þ
1=2
2 0:999822348= 0:0188549682 þ 0:9998223482 ¼ 1:879753087
k42 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2608141372 0:999822348
1=2
þ 2 ð0:018854968Þ= 0:0188549682 þ 0:9998223482 ¼ 6:395189464
t ¼ t0 þ h ¼ 1:253314137 þ 0:01 ¼ 1:263314137

x1 ¼ x10 þ hy10 þ h2 ð6k21 k31 þ 2k41 Þ=14 ¼ 0 þ 0:01 ð2:506628275Þ þ 0:012
ð1:960389968 6 þ 1:920310529 1:879753087 2Þ=14
¼ 0:025163436
x2 ¼ x20 þ hy20 þ h2 ð6k22 k32 þ 2k42 Þ=14 ¼ 1 þ 0:01 0 þ 0:012
ð6:320698363 6 þ 6:3580314546:395189464 2Þ=14
¼ 0:999683167
k51 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2633141372 ð0:025163436Þ
1=2
2 0:999683167= 0:0251634362 þ 0:9996831672 ¼ 1:838727088
k52 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2633141372 0:999683167 þ 2
1=2
ð0:025163436Þ= 0:0251634362 þ 0:9996831672 ¼ 6:432154702
Now, we are ready to compute x1, x2, y1, and y2 at t = t0 + h as follows

x1 ðt0 þ hÞ ¼ x10 þ hy10 þ h2 ð7k11 þ 24k21 þ 6k31 þ 8k41 Þ=90 þ O h7
¼ 0 þ 0:01 ð2:506628275Þ þ 0:012 ð2:0 71:960389968 24
1:920310529 61:879753087 8Þ=90 ¼ 0:025163626

x2 ðt0 þ hÞ ¼ x20 þ hy20 þ h2 ð7k12 þ 24k22 þ 6k32 þ 8k42 Þ=90 þ O h7 ¼ 1 þ 0:01 0
þ 0:012 ð6:283185304 76:320698363 246:358031454 6
6:395189464 8Þ=90 ¼ 0:999683345

y1 ðt0 þ hÞ ¼ y10 þ hð7k11 þ 32k21 þ 12k31 þ 32k41 þ 7k51 Þ=90 þ O h7
¼ 2:506628275 þ 0:01 ð2:0 71:960389968 32
1:920310529 121:879753087 32 1:838727088 7Þ=90
¼ 2:525828208

y2 ðt0 þ hÞ ¼ y20 þ hð7k12 þ 32k22 þ 12k32 þ 32k42 þ 7k52 Þ=90 þ O h7 ¼ 0 þ 0:01
ð6:283185304 76:320698363 326:358031454 12
6:395189464 326:432154702 7Þ=90 ¼ 0:063579129
These values are to be compared with those computed analytically in the same
point, which are given below:
x1 ðt0 þ hÞ ¼ 0:025163625
x2 ðt0 þ hÞ ¼ 0:999683345
y1 ðt0 þ hÞ ¼ 2:525828205
y2 ðt0 þ hÞ ¼ 0:063579126
6.6 Step-Size Control with Runge–Kutta–Nyström

Methods
In 1972, Erwin Fehlberg found Runge–Kutta–Nyström formulae with automatic

step-size control based on the leading term of the local truncation error of the
position vector only.
He found formulae for orders from four to eight. As an example, fourth-order
formulae found by Fehlberg [10] with fifth-order control on the position vector x
are shown below:
1 2
xi þ 1 ¼ xi þ h y i þ h ð13k1 þ 36k2 þ 9k3 þ 2k4 Þ þ O h5
120
1
yi þ 1 ¼ yi þ hðk1 þ 3k2 þ 3k3 þ k4 Þ þ O h5
8
1 2
x i þ 1 ¼ xi þ h y i þ h ð13k1 þ 36k2 þ 9k3 þ 2k5 Þ þ O h6
120
where x*i+1 is only computed for step-size control and

1 1 1 2
3 3 18

2 2 2
k3 ¼ f ti þ h; xi þ h yi þ h2 k2
3 3 9

1 2
k4 ¼ f ti þ h; xi þ h yi þ h ð2k1 þ k3 Þ
6

1 2
k5 ¼ f ti þ h; xi þ h yi þ h ð13k1 þ 36k2 þ 9k3 þ 2k4 Þ
120
The formulae shown above require five evaluations of the integrand function for the
first step only, because k5 is the same as k1 for every step after the first.
In 1973, Bettis [11] found formulae which use an estimate of the local truncation
error of both the position and velocity vectors to obtain step-size control. Fourth-order
formulas found by him with fifth-order control on x and y are shown below:
1
xi þ 1 ¼ xi þ h yi þ h2 ðk1 þ 2k4 Þ þ O h5
6
1
yi þ 1 ¼ yi þ hð2k3 k4 þ 2k5 Þ þ O h5
3
1 2
x i þ 1 ¼ xi þ h y i þ h ð7k1 þ 24k3 þ 6k4 þ 8k5 Þ þ O h6
90
1
y i þ 1 ¼ yi þ hð7k1 þ 32k3 þ 12k4 þ 36k5 þ 7k6 Þ þ O h6
90
6.6 Step-Size Control with Runge–Kutta–Nyström Methods 899
where the solution of higher order x*i+1 and its derivative y*i+1 are only computed
for step-size control, and

1 1 1 2
8 8 128

1 1 1 2
k3 ¼ f ti þ h; xi þ h yi þ h ðk1 þ 2k2 Þ
4 4 96

1 1 1 2
k4 ¼ f ti þ h; xi þ h yi þ h ðk1 þ 2k3 Þ
2 2 24

3 3 1 2
k5 ¼ f ti þ h; xi þ h yi þ h ð9k1 þ 18k3 þ 9k4 Þ
4 4 128

1 2
k6 ¼ f ti þ h; xi þ h yi þ h ð7k1 þ 24k3 þ 6k4 þ 8k5 Þ
90
The Runge–Kutta coefficients found by Bettis [11] are given in the following
table.
c1 = 0 c2 = 1/8 c3 = 1/4 c4 = 1/2 c5 = 3/4 c6 = 1

b^1 = 1/6 b^2 = 0 b^3 = 0 b^4 = 1/3
b1 = 0 b2 = 0 b3 = 2/3 b4 = −1/3 b5 = 2/3
b^*1 = 7/90 b^*2 = 0 b^*3 = 4/15 b^*4 = 1/15 b^*5 = 4/15
b*1 = 7/90 b*2 = 0 b*3 = 16/45 b*4 = 2/15 b*5 = 16/45 b*6 = 7/90
a11 = 1/128
a21 = 1/96 a22 = 1/48
a31 = 1/24 a32 = 0 a33 = 1/12
a41 = 9/128 a42 = 0 a43 = 9/64 a44 = 9/128
a51 = 7/90 a52 = 0 a53 = 4/15 a54 = 1/15 a55 = 4/45
where the hat sign (^) is used for the x formula, and the asterisk (*) is used for
the step-size control formulae of both x and y. Local truncation error control is
performed by evaluating the following vectors
X
s
^ i þ 1 ¼ h2 ðb^ k b^ k Þkk
k¼1
X
s
i þ 1 ¼ h ðbk b k Þkk
k¼1
As an example of application of these formulae, we consider the same system of two

second-order differential equations as that proposed by Bettis [11], in scalar form
2x2
x00 1 ¼ f1 ðt; x1 ; x2 Þ ¼ 4t2 x1 12
x21 þ x22
2x1
x00 2 ¼ f2 ðt; x1 ; x2 Þ ¼ 4t2 x2 þ 12
x21 þ x22
with the following initial conditions

p12
t0 ¼ ¼ 1:253314137
2
x10 ¼ 0:0
x20 ¼ 1:0
x0 10 y10 ¼ ð2pÞ2 ¼ 2:506628275

1
x0 20 y20 ¼ 0:0
As has been shown in the preceding paragraph, the exact solution is
x1 ðtÞ ¼ cosðt2 Þ
x2 ðtÞ ¼ sinðt2 Þ
The first step of numerical integration, with a step size h = 0.001, may be per-
formed as shown below.
t ¼ t0 ¼ 1:253314137
x1 ¼ x10 ¼ 0:0
x2 ¼ x20 ¼ 1:0
1=2
k11 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2533141372 02 1= 02 þ 12 ¼ 2:0
2
2 1=2
k12 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:253314137 1 þ 2 0= 0 þ 1
2
¼ 6:283185304
t ¼ t0 þ h=8 ¼ 1:253314137 þ 0:001=8 ¼ 1:253439137

x1 ¼ x10 þ hy10 =8 þ h2 k11 =128 ¼ 0 þ 0:001 ð2:506628275Þ=8
þ 0:0012 ð2:0Þ=128 ¼ 0:000313344
x2 ¼ x20 þ hy20 =8 þ h2 k12 =128 ¼ 1 þ 0:001 0=8 þ 0:0012
ð6:283185304Þ=128 ¼ 0:99999995
k21 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2534391372 ð0:000313344Þ
1=2
2 0:99999995= 0:0003133442 þ 0:999999952 ¼ 1:998030709
k22 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2534391372 0:99999995
1=2
þ 2 ð0:000313344Þ= 0:0003133442 þ 0:999999952 ¼ 6:285062226
t ¼ t0 þ h=4 ¼ 1:253314137 þ 0:001=4 ¼ 1:253564137

x1 ¼ x10 þ hy10 =4 þ h2 ðk11 þ 2k21 Þ=96 ¼ 0 þ 0:001 ð2:506628275Þ=4 þ 0:0012
ð2:02 1:998030709Þ=96 ¼ 0:000626719
x2 ¼ x20 þ hy20 =4 þ h2 ðk12 þ 2k22 Þ=96 ¼ 1 þ 0:001 0=4 þ 0:0012
ð6:2831853042 6:285062226Þ=96 ¼ 0:999999803
k31 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2535641372 ð0:000626719Þ
1=2
2 0:999999803= 0:0006267192 þ 0:9999998032 ¼ 1:996060245
k32 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2535641372 0:999999803
1=2
þ 2 ð0:000626719Þ= 0:0006267192 þ 0:9999998032 ¼ 6:286944382
t ¼ t0 þ h=2 ¼ 1:253314137 þ 0:001=2 ¼ 1:253814137

x1 ¼ x10 þ hy10 =2 þ h2 ðk11 þ 2k31 Þ=24 ¼ 0 þ 0:001 ð2:506628275Þ=2 þ 0:0012
ð2:0 2 1:996060245Þ=24 ¼ 0:001253563
x2 ¼ x20 þ hy20 =2 þ h2 ðk12 þ 2k32 Þ=24 ¼ 1 þ 0:001 0=2 þ 0:0012
ð6:2831853042 6:286944382Þ=24 ¼ 0:999999214
k41 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2538141372 ð0:001253563Þ
1=2
2 0:999999214= 0:0012535632 þ 0:9999992142 ¼ 1:992115774
k42 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2538141372 0:999999214
1=2
þ 2 ð0:001253563Þ= 0:0012535632 þ 0:9999992142 ¼ 6:290701744
t ¼ t0 þ 3h=4 ¼ 1:253314137 þ 3 0:001=4 ¼ 1:254064137

x1 ¼ x10 þ 3hy10 =4 þ h2 ð9k11 þ 18k31 þ 9k41 Þ=128 ¼ 0 þ 3 0:001
ð2:506628275Þ=4 þ 0:0012
ð2:0 91:996060245 181:992115774 9Þ=128
¼ 0:001880532
x2 ¼ x20 þ 3hy20 =4 þ h2 ð9k12 þ 18k32 þ 9k42 Þ=128 ¼ 1 þ 3 0:001 0=4
þ 0:0012 ð6:283185304 96:286944382 186:290701744 9Þ=128
¼ 0:99999828
k51 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2540641372 ð0:001880532Þ
1=2
2 0:99999828= 0:0018805322 þ 0:999998282 ¼ 1:988166585
k52 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2540641372 0:99999828
1=2
þ 2 ð0:001880532Þ= 0:0018805322 þ 0:999998282 ¼ 6:294457683
t ¼ t0 þ h ¼ 1:253314137 þ 0:001 ¼ 1:254314137

x1 ¼ x10 þ hy10 þ h2 ð7k11 þ 24k31 þ 6k41 þ 8k51 Þ=90 ¼ 0 þ 0:001 ð2:506628275Þ
þ 0:0012 ð2:0 71:996060245 241:992115774 6
1:988166585 8Þ=90 ¼ 0:002507625
x2 ¼ x20 þ hy20 þ h2 ð7k12 þ 24k32 þ 6k42 þ 8k52 Þ=90 ¼ 1 þ 0:001 0 þ 0:0012
ð6:283185304 7 6:286944382 24 6:290701744 6
6:294457683 8Þ=90 ¼ 0:999996855
k61 ¼ f1 ðt; x1 ; x2 Þ ¼ 4 1:2543141372 ð0:002507625Þ
1=2
2 0:999996855= 0:0025076252 þ 0:9999968552 ¼ 1:984212687
k62 ¼ f2 ðt; x1 ; x2 Þ ¼ 4 1:2543141372 0:999996855
1=2
þ 2 ð0:002507625Þ= 0:0025076552 þ 0:9999968552 ¼ 6:298211275

x1 ðt0 þ hÞ ¼ x10 þ hy10 þ h2 ð7k11 þ 2k41 Þ=6 þ O h5 ¼ 0 þ 0:001 ð2:506628275Þ
þ 0:0012 ð2:02 1:992115774Þ=6 ¼ 0:002507625

x2 ðt0 þ hÞ ¼ x20 þ hy20 þ h2 ð7k12 þ 2k42 Þ=6 þ O h5 ¼ 1 þ 0:001 0 þ 0:0012
ð6:2831853046:290701744 2Þ=6 ¼ 0:999996855

y1 ðt0 þ hÞ ¼ y10 þ hð2k31 k41 þ 2k51 Þ=3 þ O h5
¼ 2:506628275 þ 0:001 ð1:996060245 2 þ 1:992115774
1:988166583 2Þ=3 ¼ 2:508620388

y2 ðt0 þ hÞ ¼ y20 þ hð2k32 k42 þ 2k52 Þ=3 þ O h5
¼ þ 0:001 ð6:286944382 2 þ 6:2907017446:294457683 2Þ=3
¼ 0:0062907

x 1 ðt0 þ hÞ ¼ x10 þ hy10 þ h2 ð7k11 þ 24k31 þ 6k41 þ 8k51 Þ=90 þ O h6
¼ 0 þ 0:001 ð2:506628275Þ þ 0:0012 ð2:0 71:996060245 24
1:992115774 6 1:988166585 8Þ=90 ¼ 0:002507625

x 2 ðt0 þ hÞ ¼ x20 þ hy20 þ h2 ð7k12 þ 24k32 þ 6k42 þ 8k52 Þ=90 þ O h6
¼ 1 þ 0:001 0 þ 0:0012 ð6:283185304 76:286944382 24
6:290701744 66:294457683 8Þ=90 ¼ 0:999996855

y 1 ðt0 þ hÞ ¼ y10 þ hð7k11 þ 32k31 þ 12k41 þ 32k51 þ 7k61 Þ=90 þ O h6
¼ 2:506628275 þ 0:001 ð2:0 7 1:996060245 32
1:992115774 12 1:988166585 321:984212687 7Þ=90
¼ 2:508620388

y 2 ðt0 þ hÞ ¼ y20 þ hð7k12 þ 32k32 þ 12k42 þ 32k52 þ 7k62 Þ=90 þ O h6
¼ 0 þ 0:001 ð6:283185304 7 6:286944382 32
6:290701744 126:294457683 32 6:298211275 7Þ
¼ 0:0062907
These values are to be compared with those computed analytically in the same
point, which are indicated below:
x1 ðt0 þ hÞ ¼ 0:002507624
x2 ðt0 þ hÞ ¼ 0:999996855
y1 ðt0 þ hÞ ¼ 2:508620387
y2 ðt0 þ hÞ ¼ 0:006290696
6.7 Special Runge–Kutta Methods
The fundamental concepts on the methods of the Runge–Kutta family have been
given in the preceding paragraphs. The present paragraph and the following ones
show the applications of the Runge–Kutta methods to problems of astrodynamics.
The differential equations governing these problems are essentially second-order
equations
x00 ¼ f ðt; x; yÞ
in which the acceleration vector (x″) is a function of time (t), position (x), and
velocity (y), and the position and velocity vectors are functions of time. Following
the notation used in the preceding paragraphs, these equations are here written, in
the general case, in the form
x0 ¼ y
y0 ¼ f ðt; x; yÞ
with the initial conditions x(t0) = x0 and y(t0) = y0.

Here, x and y are column vectors, whose components are, respectively,
2 3 2 03
x x
6 7 6 07
x 4y5 y 4y 5
z z0
In other words, the three components of the position vector x are x, y, z and the
three components of the velocity vector y are x′, y′, z′. Hence, the two vector
equations written above correspond to six scalar equations of the first order.
Alternatively, the equations of motion are expressed in terms of the state vector z,
whose components are
23
x
6y7
6 7
6z7
z6 7
6 x0 7
6 07
4y 5
z0
According to Hairer et al. [12] and Marthinsen [13], let s be an integer called
“number of stages”, and let a21, a31, a32, …, as1, as2, as3, …, as,s−1, b1, b2, …, bs, c1,
c2, …, cs be real coefficients. An explicit Runge–Kutta method with s stages (i.e.,
s evaluations of the integrand function), for the differential equation
z0 ¼ f ðt; zÞ
with the initial condition z(t0) = z0, is defined by the following formulae:
k1 ¼ f ð t n ; z n Þ
k2 ¼ f ðtn þ c2 h; zn þ ha21 k1 Þ
k3 ¼ f ½tn þ c3 h; zn þ hða31 k1 þ a32 k2 Þ
..
.
ks ¼ f tn þ cs h; zn þ h as1 k1 þ as2 k2 þ þ as;s1 ks1
z n þ 1 ¼ z n þ hð b1 k1 þ b2 k2 þ þ b s ks Þ
A method of this type is said to be explicit, because each ki depends only on the
values ki−1, ki-2, …, k1 computed previously (this is because aij = 0 for all j
i,
which means that the matrix A {aij} of the coefficients aij is a strictly lower tri-
angular matrix), and where the ci coefficients usually satisfy the following conditions:
c2 ¼ a21
c3 ¼ a31 þ a32
..
.
cs ¼ as1 þ þ as;s1
or
X
i1
ci ¼ aij
j¼1
These coefficients are often represented by the following table (called Butcher’s
tableau, after John Butcher):
6.7 Special Runge–Kutta Methods 905
and are chosen in such a way as to make the order of the leading term of the local
truncation error of degree (p + 1) in h. They are usually determined so that
X
s
bj ¼ 1
j¼1
For the classical Runge–Kutta formulae of the fourth order, the number s of stages
is equal to the order p of the local truncation error, but such is not generally the
case. For p
5, no explicit Runge–Kutta method exists of order p with
s = p stages [12]. Butcher has shown that at least one more stage is required for
methods of order 5 and 6; that two more stages are required for methods of order 7;
and that three more stages are required for methods of order 8 and upwards. These
rules are known as Butcher’s barriers.
Of course, the values of k1, k2, …, ks must be computed at each step of the
integration process, and consequently, some authors (e.g., Sharp and Verner [14])
prefer to write the Runge–Kutta formulae as follows
!
iP
1
kni ¼ f tn þ cs h; zn þ h aij knj ði ¼ 1; 2; . . .; sÞ
j¼1
X
s
zn þ 1 ¼ zn þ h bi kni
i¼1
The methods indicated above do not possess embedded formulae for step-size
control. To this end, it is necessary to use such methods as those found by Fehlberg
or Verner (described in Sect. 6.4), or those found by Dormand and Prince [15],
which are described below.
With the embedded Runge–Kutta formulae, the same values k1, k2, …, ks are
used to compute a first approximation (of order p) to zn+1 and a second approxi-
mation (of order p + 1) to zn, the latter being obtained by means of another set of
coefficients b*1, b*2, …, b*s. The difference is in the fact that Fehlberg uses the
formula of order p to compute zn+1 and checks the results against the corresponding
value coming from the formula of order p + 1, whereas Dormand and Prince do the
reverse. Be that as it may, the notation b1, b2, …, bs is used here to indicate the
coefficients employed to advance the solution from tn to tn+1, whereas the notation
b*1, b*2, …, b*s indicates the coefficients used only for step-size control. The
Runge–Kutta methods based on embedded p(p + 1)s formulae (or embedded p + 1
(p)s formulae) are expressed by
z½p n þ 1 z½p ðtn þ hÞ ¼ zn þ hðb 1 k1 þ b 2 k2 þ . . . þ b s ks Þ

z½p þ 1 n þ 1 z½p þ 1 ðtn þ hÞ ¼ zn þ hðb1 k1 þ b2 k2 þ . . . þ bs ks Þ
where the superscripts [p] and [p + 1] indicate, respectively, the approximation of

order p and the approximation of order p + 1, s indicates the number of stages, and
k 1 ¼ f ðt n ; zn Þ
..
.
ks ¼ f tn þ cs h; zn þ h as1 k1 þ as2 k2 þ þ as;s1 ks1
are the stages or function evaluations.

The difference

½p þ 1

z ½p
nþ1 z nþ1

provides an approximation for the leading term of the local truncation error relating
to the formula of order p. This approximation is used to estimate the subsequent
step size, in order to keep the local truncation error under control.
The coefficients {bi}, {b*i}, {ci}, and {aij}, where i, j = 1, 2,…, s, are deter-
mined by imposing conditions which result from Taylor-series expansions.
For example, RK5(4)7FM by Dormand and Prince [15] is a seven-stage method
of the fifth order with an embedded formula of the fourth order for truncation error
control, based on the following formulae and coefficients:
k1 ¼ f ðtn ; zn Þ
..
.
k7 ¼ f ½tn þ c7 h; zn þ hða71 k1 þ a72 k2 þ þ a76 k6 Þ
z½5 n þ 1 ¼ zn þ hðb1 k1 þ b2 k2 þ þ b7 k7 Þ
z½4 n þ 1 ¼ zn þ hðb 1 k1 þ b 2 k2 þ þ b 7 k7 Þ
The difference

½5

z n þ 1 z½4 n þ 1
provides the leading term of the local truncation error.

The formulae to be used for computing k1, k2, …, k7, zn þ 1 ½5 , and zn þ 1 ½4 by
means of the RK5(4)7FM method are also given in Sect. 6.9.
The Butcher tableau for RK5(4)7FM by Dormand and Prince [15] is given
below.
Coefficients for the embedded Dormand and Prince RK5(4)7FM method are
–
– –
– –
–
–
–
Another example is RK8(7)13M by Prince and Dormand, which is a

thirteen-stage method of the eighth order with an embedded formula of the seventh
order for truncation error control, based on the following coefficients [16]:
Coefficients for the embedded Prince and Dormand RK8(7)13M method
c1 ¼ 0
c2 ¼ 1=18
c3 ¼ 1=12
c4 ¼ 1=8
c5 ¼ 5=16
c6 ¼ 3=8
c7 ¼ 59=400
c8 ¼ 93=200
c9 ¼ 5490023248=9719169821
c10 ¼ 13=20
c11 ¼ 1201146811=1299019798
c12 ¼ 1
c13 ¼ 1
a21 ¼ c2 ¼ 1=18
a31 ¼ 1=48
a32 ¼ 1=16
a41 ¼ 1=32
a42 ¼ 0:0
a43 ¼ 3=32
a51 ¼ 5=16
a52 ¼0
a53 ¼ 75=64
a54 ¼ a53 ¼ 75=64
a61 ¼ 3=80
a62 ¼0
a63 ¼0
a64 ¼ 3=16
a65 ¼ 3=20
a71 ¼ 29443841=614563906
a72 ¼0
a72 ¼0
a74 ¼ 77736538=692538347
a75 ¼ 28693883=1125:D6
a76 ¼ 23124283=18:D8
a81 ¼ 16016141=946692911
a82 ¼ 0
a83 ¼ 0
a84 ¼ 61564180=158732637
a85 ¼ 22789713=633445777
a86 ¼ 545815736=2771057229
a87 ¼ 180193667=1043307555
a91 ¼ 39632708=573591083
a92 ¼ 0
a93 ¼ 0
a94 ¼ 433636366=683701615
a95 ¼ 421739975=2616292301
a96 ¼ 100302831=723423;059
a97 ¼ 790204164=839813087
a98 ¼ 800635310=3783071287
a101 ¼ 246121993=1340847787
a102 ¼ 0
a103 ¼ 0
a10;4 ¼ 37695042795=15268766246
a105 ¼ 309121744=1061227803
a106 ¼ 12992083=490766935
a107 ¼ 6005943493=2108947869
a108 ¼ 393006217=1396673457
a109 ¼ 123872331=1001029789
a111 ¼ 1028468189=846180014
a112 ¼ 0
a113 ¼ 0
a114 ¼ 8478235783=508512852
a115 ¼ 1311729495=1432422823
a116 ¼ 10304129995=1701304382
a117 ¼ 48777925059=3047939560
a118 ¼ 15336726248=1032824649
a119 ¼ 45442868181=3398467696
a1110 ¼ 3065993473=597172653
a121 ¼ 185892177=718116043
a122 ¼ 0
a123 ¼ 0
a124 ¼ 3185094517=667107341
a125 ¼ 477755414=1098053517
a126 ¼ 703635378=230739211
a127 ¼ 5731566787=1027545527
a128 ¼ 5232866602=850066563
a129 ¼ 4093664535=808688257
a1210 ¼ 3962137247=1805957418
a1211 ¼ 65686358=487910083
a131 ¼ 403863854=491063109
a132 ¼ 0
a133 ¼ 0
a134 ¼ 5068492393=434740067
a135 ¼ 411421997=543043805
a136 ¼ 652783627=914296604
a137 ¼ 11173962825=925320556
a138 ¼ 13158990841=6184727034
a139 ¼ 3936647629=1978049680
a1310 ¼ 160528059=685178525
b1 ¼ 13451932=455176623
b2 ¼ 0
b3 ¼ 0
b4 ¼ 0
b5 ¼ 0
b6 ¼ 808719846=976000145
b7 ¼ 1757004468=5645159321
b8 ¼ 656045339=265891186
b9 ¼ 3867574721=1518517206
b10 ¼ 465885868=322736535
b11 ¼ 53011238=667516719
b12 ¼ 2=45
b 1 ¼ 14005451=335480064
b 2 ¼ 0
b 3 ¼ 0
b 4 ¼ 0
b 5 ¼ 0
b 6 ¼ 59238493=1068277825
b 7 ¼ 181606767=758867731
b 8 ¼ 561292985=797845732
b 9 ¼ 1041891430=1371343529
b 10 ¼ 760417239=1151165299
b 11 ¼ 118820643=751138087
b 12 ¼ 528747749=2220607170
b 13 ¼ 1=4
The approximations of, respectively, seventh and eighth order are computed by
means of
z½7 n þ 1 ¼ zn þ hðb 1 k1 þ b 2 k2 þ þ b 13 k13 Þ

z½8 n þ 1 ¼ zn þ hðb1 k1 þ b2 k2 þ þ b13 k13 Þ

The difference
zn þ 1 ½8 zn þ 1 ½7
provides the leading term of the local trun-
cation error.
6.8 Special Runge–Kutta–Nyström Methods
Let us write again the second-order differential equation x″ = y′= f(t, x, y), with the
initial conditions x(t0) = x0 and y(t0) = y0, as a system of first-order equations
6.8 Special Runge–Kutta–Nyström Methods 911
x0 ¼ y
y0 ¼ f ðt; x; yÞ
Bearing in mind that
X
i1
ci ¼ aij
j¼1
we have (see [12]):
X
s
xn þ 1 ¼ xn þ h y n þ h2 b^ i ki
i¼1
X
s
yn þ 1 ¼ yn þ h bi k i
i¼1
where
!
P
s P
s
ki ¼ f tn þ ci h; xn þ ci h yn þ h 2
a^ ij kj ; yn þ h aij kj ði ¼ 1; 2; . . .; sÞ
j¼1 j¼1
P
s P
s
a^ ij ¼ aik akj b^ i ¼ bj aji
k¼1 j¼1
The method shown above is called here, in accordance with Marthinsen [13], the
general Nyström method (or the Runge–Kutta–Nyström method, because it derives
from the Runge–Kutta method).
For example, a Butcher tableau representing the coefficients is shown below for
a general Nyström method of the fourth order.
If the right-hand side of the differential equation x″ = f(t, x, y) does not depend
explicitly on the velocity vector (y) of the body, that is, if x″ = f(t, x), then the
general Nyström method simplifies, so that the coefficients aij are not necessary any
more. Such is the case in the absence of perturbations due to atmospheric drag. In
this case, the formulae written above are modified by writing aij instead of aîj in
order to simplify the notation, as shown below.
X
s
xn þ 1 ¼ xn þ h y n þ h2 b^ i ki
i¼1
X
s
yn þ 1 ¼ yn þ h bi ki
i¼1
where
!
P
s
aij kj ði ¼ 1; 2; . . .; sÞ
j¼1
and this method is called the special Nyström method. The Butcher tableau for a
special Nyström method simplifies as shown below. Some authors, Dormand [15]
for one, use other schemes (modified Butcher tableau) for the Butcher tableau. The
scheme indicated below will be used in the present book.
For the special Nyström method, Battin [10] gives two sets of fifth-order coef-
ficients found by Nyström himself. They are shown below.
Special Nyström method of the fifth order (first set)
–
–
Special Nyström method of the fifth order (second set)
–
–
The fifth-order methods shown above require only four evaluations of the integrand
function, that is, s = 4.
The basic concepts on the special Nyström method with step-size control have
been given in Sect. 6.5. This method uses the following formulae [17]:
!
P
i1
aij kj ði ¼ 0; 1; 2; . . .; sÞ
j¼0
P
s
x n þ 1 ¼ x n þ h y n þ h 2 b i ki kk¼ Oðhp þ 1 Þ
i¼0
P
s
xn þ 1 ¼ xn þ h y n þ h2 b^ i ki kk ¼ Oðhp þ 2 Þ
i¼0
X
s
yn þ 1 ¼ yn þ h bi ki
i¼0
and the following coefficients:
where b*s = k > 0 is a free parameter which is not determined by the order
conditions.
The formulae for xn+1 and yn+1 give a (p + 1)th-order approximation of,
respectively, the position vector and the velocity vector at t = tn + h of order p + 1.
The formula for x*n+1 gives an estimate of the local truncation error vector
e jx x j ¼ kh2 jks1 ks j
and is only used to control the step size by means of the local truncation error. In
accordance with Dormand [15], the step-size control is performed by means of the
following formula:
p þ1 1
etol
hn þ 1 ¼ 0:9hn
kk
where 0.9 is a safety factor, etol is the tolerated error bound for step n + 1, and ||e|| is
some norm of the vector e indicated above. Two norms commonly used are:
(a) the Euclidean length of e
1
kk ¼ 21 þ 22 þ þ 2n 2
where e1, e2, …, en are the components of e; or

(b) the maximum in magnitude among the components of e
kk ¼ Maxi ðji jÞ

The coefficient bs = k > 0 is a free parameter which neither depends on the order
conditions nor has any effect on the solution (position and velocity vectors),
because it only appears in the step-size control formula. By changing the value of k,
it is possible to adjust the value of etol necessary to obtain selected values of
step-size and error, as will be shown below.
The conditions cs = 1 and b*i = asi (i = 0, 1, 2, …, s − 1) make it possible to
spare one evaluation of the integrand function at each step. This is because the last
evaluation of the integrand function, that is,
!
X
s1
ks ¼ f tn þ cs h; xn þ cs h yn þ h 2
asj kj
j¼0
is the same as the first evaluation

!
X
s1
k0 ¼ f ½tn þ h; xðtn þ hÞ ¼ f tn þ h; xn þ h yn þ h2 bj kj
j¼0
of the next step (which starts at tn + h). Hence, this method is said to require
s stages, even if any single step actually requires s + 1 function evaluations.
A first example of application, due to Dormand and Prince [18], is the following
table, relating to RKN7(6)T, which is a special Nyström method of the seventh
order.
Coefficients for RKN7(6)T special Nyström method (Dormand and Prince) are
c1 ¼ 1=10
c2 ¼ 1=5
c3 ¼ 3=8
c4 ¼ 1=2
c5 ¼ ½7 ð21Þ1=2 =14
c6 ¼ ½7 þ ð21Þ1=2 =14
c7 ¼1
c8 ¼1
a10 ¼ 1=200
a20 ¼ 1=150
a21 ¼ 1=75
a30 ¼ 171=8192
a31 ¼ 45=4096
a32 ¼ 315=8192
a40 ¼ 5=288
a41 ¼ 25=528
a42 ¼ 25=672
a43 ¼ 16=693
a50 ¼ ½1003 205ð21Þ1=2 =12348
a51 ¼ 25½751 173ð21Þ1=2 =90552
a52 ¼ 25½624 137ð21Þ1=2 =43218
a53 ¼ 128½361 79ð21Þ1=2 =237699
a54 ¼ ½3411 745ð21Þ1=2 =24696
a60 ¼ ½793 þ 187ð21Þ1=2 =12348
a61 ¼ 25½331 þ 113ð21Þ1=2 =90552
a62 ¼ 25½1044 þ 247ð21Þ1=2 =43218
a63 ¼ 128½14885 þ 3779ð21Þ1=2 =9745659
a64 ¼ ½3327 þ 797ð21Þ1=2 =24696
a65 ¼ ½581 þ 127ð21Þ1=2 =1722
a70 ¼ ½157 3ð21Þ1=2 =378
a71 ¼ 25½143 10ð21Þ1=2 =2772
a72 ¼ 25½876 þ 55ð21Þ1=2 =3969
a73 ¼ 1280½913 þ 18ð21Þ1=2 =596673
a74 ¼ ½1353 þ 26ð21Þ1=2 =2268
a75 ¼ 7½1777 þ 377ð21Þ1=2 =4428
a76 ¼ 7½5 ð21Þ1=2 =36
a80 ¼ 1=20
a81 ¼0
a82 ¼0
a83 ¼0
a84 ¼ 8=45
a85 ¼ 7½7 þ ð21Þ1=2 =360
a86 ¼ 7½7 ð21Þ1=2 =360
a87 ¼0
b 0 ¼ a80 ¼ 1=20
b 1 ¼ a81 ¼ 0
b 2 ¼ a82 ¼ 0
b 3 ¼ a83 ¼ 0
b 4 ¼ a84 ¼ 8=45
b 5 ¼ a85 ¼ 7½7 þ ð21Þ1=2 =360
b 6 ¼ a86 ¼ 7½7 ð21Þ1=2 =360
b 7 ¼ k
b 8 ¼k
b^ 2 ¼ b2 ¼ 0
b^ 3 ¼ b3 ¼ 0
b^ 4 ¼ b4 ¼ 8=45
b^ 5 ¼ b5 ¼ 7½7 þ ð21Þ1=2 =360
b^ 7 ¼ 0
b^ 8 ¼ 0
b0 ¼ 1=20
b1 ¼0
b2 ¼0
b3 ¼0
b4 ¼ 16=45
b5 ¼ 49=180
b6 ¼ 49=180
b7 ¼ 1=20
b8 ¼0
A second example, due to Filippi and Gräf [19], is given below. It is an

eleventh-order, seventeen-stage special Nyström method, called a RKN11(10)17
formula pair, because the solution is computed by means of an eleventh-order
formula, with a tenth-order formula used for error control and seventeen evaluations
of the integrand function. The method is based on the following coefficients.
Coefficients for RKN11(10)17 special Nyström method (Filippi and Gräf) are
c1 ¼ 0:4166666666666666666666666667D01
c2 ¼ 0:8333333333333333333333333333D01
c3 ¼ 0:2000000000000000000000000000Dþ00
c4 ¼ 0:2500000000000000000000000000Dþ00
c5 ¼ 0:7619047619047619047619047647Dþ00
c6 ¼ 0:3333333333333333333333333333Dþ00
c7 ¼ 0:1333333333333333333333333337Dþ00
c8 ¼ 0:5000000000000000000000000000Dþ00
c9 ¼ 0:8333333333333333333333333333Dþ00
c10 ¼ 0:8333333333333333333333333333D01
c11 ¼ 0:2500000000000000000000000000Dþ00
c12 ¼ 0:3333333333333333333333333333Dþ00
c13 ¼ 0:6666666666666666666666666666Dþ00
c14 ¼ 0:7500000000000000000000000000Dþ00
c15 ¼ 0:9242424242424242424242460752Dþ00
c16 ¼ 0:1000000000000000000000000000Dþ01
c17 ¼ 0:1000000000000000000000000000Dþ01
a1;0 ¼ 0:8680555555555555555555555555D03
a2;0 ¼ 0:1157407407407407407407407407D02
a2;1 ¼ 0:2314814814814814814814814815D02
a3;0 ¼ 0:1040000000000000000000000000D01
a3;1 ¼ 0:1280000000000000000000000000D01
a3;2 ¼ 0:2240000000000000000000000000D01
a4;0 ¼ 0:6510416666666666666666666668D02
a4;1 ¼ 0:0
a4;2 ¼ 0:2008928571428571428571428571D01
a4;3 ¼ 0:4650297619047619047619047619D02
a5;0 ¼ 0:7442558350050598650718971132Dþ00
a5;1 ¼ 0:0
a5;2 ¼ 0:2398395982021852194994338203Dþ01
a5;3 ¼ 0:4296409767675047353488214866Dþ01
a5;4 ¼ 0:2932519053764830987284594640Dþ01
a6;0 ¼ 0:8320473251028806584362139914D02
a6;1 ¼ 0:0
a6;2 ¼ 0:3148612271419288963148612272D01
a6;3 ¼ 0:7784553453103358941400372646D02
a6;4 ¼ 0:7962484448272561967652406931D02
a6;5 ¼ 0:1921688957938430654513332561D05
a7;0 ¼ 0:3177389300411522633744855977D02
a7;1 ¼ 0:0
a7;2 ¼ 0:7172885695720783440081685729D02
a7;3 ¼ 0:4098087864566207315736505945D02
a7;4 ¼ 0:3150419446837017896449420999D02
a7;5 ¼ 0:1337454518694045049020649864D05
a7;6 ¼ 0:5150551440329218106995884776D03
a8;0 ¼ 0:1925455729166666666666666669D01
a8;1 ¼ 0:0
a8;2 ¼ 0:0
a8;3 ¼ 0:0
a8;4 ¼ 0:1683277962347729789590254748D01
a8;5 ¼ 0:8261897683227625088090204146D04
a8;6 ¼ 0:4322916666666666666666666683D01
a8;7 ¼ 0:7926643668831168831168831192D01
a9;0 ¼ 0:3890677244084362139917695514D01
a9;1 ¼ 0:0
a9;2 ¼ 0:0
a9;3 ¼ 0:0
a9;4 ¼ 0:1162107105259560039921797094Dþ00
a9;5 ¼ 0:1280779593149063653526890507D01
a9;6 ¼ 0:2813143004115226337448560306D01
a9;7 ¼ 0:1000838847074642529187983710Dþ00
a9;8 ¼ 0:1073444886576199707512838846Dþ00
a10;0 ¼ 0:2298638614607445987654320989D02
a10;1 ¼ 0:0
a10;2 ¼ 0:0
a10;3 ¼ 0:0
a10;4 ¼ 0:0
a10;5 ¼ 0:2512712158652423387447237099D03
a10;6 ¼ 0:9907510109453589391860996819D03
a10;7 ¼ 0:1723141771801382228570780760D02
a10;8 ¼ 0:5659797947913478216508519521D03
a10;9 ¼ 0:1264842678326474622770919127D03
a11;0 ¼ 0:6848028273809523809523809520D02
a11;1 ¼ 0:0
a11;2 ¼ 0:0
a11;3 ¼ 0:0
a11;4 ¼ 0:0
a11;5 ¼ 0:0
a11;6 ¼ 0:1209077380952380952380952379D02
a11;7 ¼ 0:6982995613017934446505875037D02
a11;8 ¼ 0:2051204004329004329004328998D03
a11;9 ¼ 0:6288973922902494331065759621D05
a11;10 ¼ 0:1640873015873015873015873020D01
a12;1 ¼ 0:0
a12;2 ¼ 0:0
a12;3 ¼ 0:0
a12;4 ¼ 0:1224804770337618196047355743D01
a12;5 ¼ 0:2077286426548597692800167757D02
a12;6 ¼ 0:3084352111917326351227812139D01
a12;7 ¼ 0:1322398920232751620391476668Dþ00
a12;8 ¼ 0:5694717797969194984675336503D03
a12;9 ¼ 0:1144721852146542255214994914D02
a12;10 ¼ 0:1366017573669854049407166350Dþ00
a12;11 ¼ 0:7718593382280344668606066469D01
a13;0 ¼ 0:3147136427366806555520298306D01
a13;1 ¼ 0:0
a13;2 ¼ 0:0
a13;3 ¼ 0:0
a13;4 ¼ 0:1569719156364494947817982402Dþ00
a13;5 ¼ 0:6654987695858527013495559480D02
a13;6 ¼ 0:0
a13;7 ¼ 0:8269612470709131479069575139D01
a13;8 ¼ 0:5277046384601007307143946853D01
a13;9 ¼ 0:3564062192959960179721795389D02
a13;10 ¼ 0:3782761972062383585113875547D02
a13;11 ¼ 0:8098383997314998286826664921D01
a13;12 ¼ 0:1383011879288579268975993216D01
a14;0 ¼ 0:8551325588087360429217853897D02
a14;1 ¼ 0:0
a14;2 ¼ 0:0
a14;3 ¼ 0:0
a14;4 ¼ 0:8401264079085973583990782598D01
a14;5 ¼ 0:1186337961569032036179412516D02
a14;6 ¼ 0:1135641886457673424136874024D þ 00
a14;7 ¼ 0:2190092537441172654013264721Dþ00
a14;8 ¼ 0:4620369354790417202229644247D01
a14;9 ¼ 0:1120372383953845501878892251D02
a14;10 ¼ 0:2618995077705879099608244632Dþ00
a14;11 ¼ 0:2841919156005766259624988398Dþ00
a14;12 ¼ 0:1007728344392311617920989094Dþ00
a14;13 ¼ 0:1325349183291664781212037926D01
a15;0 ¼ 0:3227323796009178105251112707D01
a15;1 ¼ 0:0
a15;2 ¼ 0:0
a15;3 ¼ 0:0
a15;4 ¼ 0:6552253393177109965556635085D03
a15;5 ¼ 0:2053435081848825845464654784D02
a15;6 ¼ 0:8183285936474925294914041188D01
a15;7 ¼ 0:9687035581635444102020453817D01
a15;8 ¼ 0:1009741982188193601657085464Dþ00
a15;9 ¼ 0:7943878492642055987892133169D02
a15;10 ¼ 0:3988879155234310326668935987D01
a15;11 ¼ 0:6260072946369460370714950309D01
a15;12 ¼ 0:2982691974303517060774392676D01
a15;13 ¼ 0:1160314062805673150276068512D01
a15;14 ¼ 0:2566041805220703528164152583D01
a16;0 ¼ 0:1598229814960049517404649275D 01
a16;1 ¼ 0:0
a16;2 ¼ 0:0
a16;3 ¼ 0:0
a16;4 ¼ 0:0
a16;5 ¼ 0:0
a16;6 ¼ 0:0
a16;7 ¼ 0:0
a16;8 ¼ 0:5685521470453431953596264635D01
a16;9 ¼ 0:2342953776266563711445070295D01
a16;10 ¼ 0:1485105390278212809150227904Dþ00
a16;11 ¼ 0:3889673830278987492323711070D01
a16;12 ¼ 0:1485093897141690039134706881Dþ00
a16;13 ¼ 0:3093315321482061128468761163D01
a16;14 ¼ 0:6806447423483744354174394038D01
a16;15 ¼ 0:1567773041409260782627942270D01
a17;0 ¼ 0:2352719229768410096278945558D01
a17;1 ¼ 0:0
a17;2 ¼ 0:0
a17;3 ¼ 0:0
a17;4 ¼ 0:0
a17;5 ¼ 0:0
a17;6 ¼ 0:0
a17;7 ¼ 0:0
a17;8 ¼ 0:1095238095238095238095226633Dþ00
a17;9 ¼ 0:5714285714285714285717006159D02
a17;10 ¼ 0:1235521235521235521235522436Dþ00
a17;11 ¼ 0:1287533440342429106474043912Dþ00
a17;12 ¼ 0:3989010989010989010989112569D01
a17;13 ¼ 0:2382352941176470588235556148D01
a17;14 ¼ 0:3621808143547273982056220803D01
a17;15 ¼ 0:8997524140506862358205344942D02
a17;16 ¼ 0:0
b 0 ¼ a17;0 ¼ 0:2352719229768410096278945558D01
b 1 ¼ a17; 1 ¼ 0:0
b 2 ¼ a17; 2 ¼ 0:0
b 3 ¼ a17; 3 ¼ 0:0
b 4 ¼ a17; 4 ¼ 0:0
b 5 ¼ a17; 5 ¼ 0:0
b 6 ¼ a17; 6 ¼ 0:0
b 7 ¼ a17; 7 ¼ 0:0
b 8 ¼ a17; 8 ¼ 0:1095238095238095238095226633Dþ00
b 9 ¼ a17; 9 ¼ 0:5714285714285714285717006159D02
b 10 ¼ a17; 10 ¼ 0:1235521235521235521235522436Dþ00
b 11 ¼ a17; 11 ¼ 0:1287533440342429106474043912Dþ00
b 12 ¼ a17; 12 ¼ 0:3989010989010989010989112569D01
b 13 ¼ a17; 13 ¼ 0:2382352941176470588235556148D01
b 14 ¼ a17; 14 ¼ 0:3621808143547273982056220803D01
b 15 ¼ a17; 15 ¼ 0:8997524140506862358205344942D02
b 16 ¼ k
b 17 ¼ k
b^ 0 ¼ b0 ¼ 0:2352719229768410096278945558D 01
b^ 1 ¼ b1 ¼ 0:0
b^ 2 ¼ b2 ¼ 0:0
b^ 3 ¼ b3 ¼ 0:0
b^ 4 ¼ b4 ¼ 0:0
b^ 5 ¼ b5 ¼ 0:0
b^ 7 ¼ b7 ¼ 0:0
b^ 8 ¼ b8 ¼ 0:01095238095238095238095226633Dþ00
b^ 9 ¼ b9 ¼ 0:5714285714285714285717006159D02
b^ 10 ¼ b10 ¼ 0:1235521235521235521235522436Dþ00
b^ 11 ¼ b11 ¼ 0:1287533440342429106474043912Dþ00
b^ 12 ¼ b12 ¼ 0:3989010989010989010989112569D01
b^ 13 ¼ b13 ¼ 0:2382352941176470588235556148D01
b^ 14 ¼ b14 ¼ 0:3621808143547273982056220803D01
b^ 15 ¼ b15 ¼ 0:8997524140506862358205344942D02
b^ 16 ¼ 0:0
b^ 17 ¼ 0:0
b0 ¼ 0:2352719229768410096278945558D 01
b1 ¼ 0:0
b2 ¼ 0:0
b3 ¼ 0:0
b4 ¼ 0:0
b5 ¼ 0:0
b6 ¼ 0:0
b7 ¼ 0:0
b8 ¼ 0:2190476190476190476190453266Dþ00
b9 ¼ 0:3428571428571428571430203695D01
b10 ¼ 0:1347841347841347841347842657Dþ00
b11 ¼ 0:1716711253789905475298725216Dþ00
b12 ¼ 0:5983516483516483516483668854D01
b13 ¼ 0:7147058823529411764706668444D01
b14 ¼ 0:1448723257418909592822488321Dþ00
b15 ¼ 0:1187673186546905831283162769Dþ00
b16 ¼ 0:2173881673881673881673791154D01
b17 ¼ 0:0
where the letter D preceded by a floating-point number and followed by a positive

or negative integer means that the number preceding D is to be multiplied by ten
raised to the power specified by the integer which follows D (e.g., c3 = 0.2D+00
means 0.2 100 = 0.2).
For this method, Filippi and Gräf give an estimate of the local truncation error in
accordance with the formula shown above relating to the general case of a method
with s stages. By taking s = 17, we have
e jx x j ¼ kh2 jk16 k17 j
As indicated by Filippi and Gräf [19], the value of k is more or less arbitrary.
With a tolerance of 10–24 to achieve a global error of about 10–20, they suggest to
choose k = 1.4 10−5 for the set of coefficients given above.
6.9 Interpolants
Ordinary explicit Nyström methods are usually applied to approximate the solution
of a second-order differential equation x″= f(t, x, y) on a mesh of selected points.
The interval [t0, tF] of integration comprises a mesh of points t0, t1, t2, …, tF, where
the distance between each point tn and the contiguous point tn+1 depends on the
method used for step-size selection. The approximate solution xn x(tn) is com-
puted at the mesh points, and the steps are generally taken as large as possible,
provided that the requirements of stability and accuracy are met.
However, sometimes the solution is required at points which do not belong to the
mesh, for example, at a certain time s which is contained in an interval [tn, tn+1],
where tn and tn+1 are points of the mesh. In the problems of orbit computations
considered here, such cases arise when the data of a satellite orbit must be repre-
sented graphically, or when an iterative search for special events (e.g., the points at
which a satellite enters into or exits from the shadow of the Earth) must be per-
formed. The simplest way to solve the problem is to take, starting from tn, a step of
length s − tn smaller than tn+1 − tn, so as to reach just the desired time s. However,
sometimes the particular times t are not known in advance. Another possible
solution is to integrate the equation at the mesh points and to interpolate the results
by means of appropriate polynomials. However, this implies the necessity of storing
the results for the subsequent interpolation. These facts justify the search for
effective methods to compute the results at points arbitrarily chosen within a given
interval. Following Marthinsen [13], we call these methods interpolants; other
authors (Dormand, for one) call them dense output. The problem of interpolants has
been studied by various authors, among whom are Horn, Dormand, Enright,
Marthinsen, Montenbruck and Gill, Owren and Zennaro, …, whose papers are cited
at the end of this chapter.
6.9 Interpolants 925
The fundamental idea is to use the integrand function evaluations (k1, k2, …, ks)
and a few additional evaluations (ks+1, ks+2, …) to calculate an interpolating
polynomial. This polynomial is used to compute the solution at any point t placed
between two contiguous mesh points tn and tn+1.
One of the earliest continuous methods has been proposed by Horn in 1981 [20]
and is based on the embedded six-stage Runge–Kutta–Fehlberg method RK4(5)6
[4]. The formulae and coefficients of RK4(5)6 have been given, in scalar form, in
Sect. 6.3. They are also given below in vector form for convenience of the reader.
k1 ¼f ðtn ;xn Þ

1 1
4 4

3 3 9
k3 ¼ f tn þ h; xn þ h k1 þ h k2
8 32 32

12 1932 7200 7296
k4 ¼ f tn þ h; xn þ h k1 h k2 þ h k3
13 2197 2197 2197

439 3680 845
k5 ¼ f tn þ h; xn þ h k1 8h k2 þ h k3 h k4
216 513 4104

1 8 3544 1859 11
k6 ¼ f tn þ h; xn h k1 þ 2h k2 h k3 þ h k4 h k5
2 27 2565 4104 40

25 1408 2197 1
xn þ 1 ½4 ¼ xn þ h k1 þ k3 þ k4 k 5
216 2565 4104 5

½5 16 6656 28561 9 2
xn þ 1 ¼ x n þ h k1 þ k3 þ k4 k5 þ k6
135 12825 56430 50 55
Let us write the differential equation x″ = f(t, x, y) in the form z′ = f(t, z), where
z is the state vector. The fourth-order solution and the fifth-order solution are
P
5 P
6
zn þ 1 ½4 ¼ zn þ h bi ki zn þ 1 ½5 ¼ zn þ h b i ki
i¼1 i¼1
where
!
iP
1
ki ¼ f tn þ ci h; zn þ h aij kj ði ¼ 1; 2; . . .; 6Þ
j¼1
The Horn’s continuous method is based on the RK4(5)6 formulae. This method
performs one more stage than the original (i.e., discrete) RK4(5)6, that is, one more
evaluation (k7) of the integrand function, for the purpose of computing the solution
at any time t between tn and tn + h. The additional stage k7 is

1 1 2
k7 ¼ f tn þ h; zn þ h k1 þ k5 þ k6
6 6 3
After the computation of the six regular stages k1, k2, …, k6 and of the additional
stage k7, let us suppose that the approximate solution z[4] must be computed at a
point s = tn + rh, where r is a number such that 0 < r < 1. This approximation is
X
7
z½4 ðtÞ ¼ z½4 ðtn þ rhÞ ¼ zn þ rh bi ki
i¼1
where the continuous weights bi = bi(r) are polynomials in the independent vari-
able r 2 [0, 1]. The expression of the continuous weights (from [20]) as a function
of r is

301 269 311
b1 ¼ 1 r þr þ r
120 108 360
b2 ¼ 0

7168 4096 14848
b3 ¼ r þr þ r
1425 513 4275

28561 199927 371293
b4 ¼r þr r
8360 22572 75240

57 42
b5 ¼r þ r 3 þ r
50 25

96 40 102
b6 ¼r þr r
55 11 55

3 5
b7 ¼ r þ r 4 þ r
2 2
Another method is based on Dormand-Prince’s RK5(4)7FM (see Sect. 6.7),

instead of Fehlberg’s RK4(5)6. The coefficients, which have been given in Sect. 6.7
in form of Butcher’s tableau, concern the discrete fifth-order method.
The formulae based on these coefficients are also given below for convenience.
k1 ¼ f ðtn ; zn Þ

1 1
k2 ¼ f tn þ h; zn þ h k1
5 5

3 3 9
k3 ¼ f tn þ h; zn þ h k1 þ h k2
10 40 40

4 44 56 32
k4 ¼ f tn þ h; zn þ h k1 h k2 þ hk3
5 45 15 9

8 19372 25360 6448 212
k5 ¼ f tn þ h; zn þ h k1 h k2 þ h k3 h k4
9 6561 2187 6561 729

9017 355 46732 49 5103
k6 ¼ f tn þ h; zn þ h k1 h k2 þ h k3 þ h k4 h k5
3168 33 5247 176 18656

35 500 125 2187 11
k7 ¼ f tn þ h; zn þ h k1 þ h k3 þ h k4 h k5 þ h k6
384 1113 192 6784 84

35 500 125 2187 11
zn þ 1 ½4 ¼ zn þ h k1 þ k3 þ k4 k5 þ k6
384 1113 192 6784 84

5179 7571 393 92097 187 1
zn þ 1 ½5 ¼ zn þ h k1 þ k3 þ k4 k5 þ k6 þ k7
57600 16695 640 339200 2100 40
Owren and Zennaro [21] have considered a fourth-order continuous extension to

RK5(4)7FM, obtained by using the first six stages of RK5(4)7FM. This continuous
extension is
z½4 ðsÞ ¼ z½4 ðtn þ rhÞ ¼ zn þ hðb1 k1 þ b2 k2 þ þ b6 k6 Þ
where the continuous weights b1, b2, …, b6 are polynomials in r and are given
below (from [21]).
1337 2 1039 3 1163 4

b1 ¼ r r þ r r
480 360 1152
b2 ¼ 0
4216 2 18728 3 7580 4
b3 ¼ r r þ r
1113 3339 3339
27 9 415 4
b4 ¼ r2 þ r3 r
16 2 192
2187 2 2763 3 8991 4
b5 ¼ r þ r r
8480 2120 6784
33 319 3 187 4
b6 ¼ r2 r þ r
35 105 84
Starting from these results, Owren has found [22, 23] a nine-stage, fifth-order
continuous extension of RK5(4)7FM. The Butcher tableau is given below.
Owren’s fifth-order continuous extension to RK5(4)7FM
–
– –
– –
–
– –
–
–
–
The fifth-order continuous weights for RK5(4)7FM computed by Owren are
285 2 97 3 813 4 29 5
b1 ¼ r r þ r r þ r
64 12 128 16
b2 ¼ 0
1000 2 16000 3 8500 4 4000 5
b3 ¼ r r þ r r
371 1113 371 371
125 2 125 3 2125 4 125 5
b4 ¼ r r þ r r
32 6 64 8
6561 2 2187 3 111537 4 6561 5
b5 ¼ r þ r r þ r
3392 212 6784 848
11 88 187 4 22 5
b6 ¼ r2 r3 þ r r
14 21 28 7
7 2 19 3 63 4
b7 ¼ r þ r r þ 4r5
8 4 8
125 2 125 3 125 4
b8 ¼ r r þ r
24 12 24
16 80 3 112 4
b9 ¼ r2 þ r r þ 16r5
3 3 3
Dormand, too, has found continuous extensions [24] to his discrete RK5(4)7FM
formula pair, but there is a difference between his version and that found by Owren.
This difference is in the values of the a8j and a9j coefficients, relating to the two
additional stages. As Owren [23] points out, the continuous fifth-order weights are
equal in both versions, but Dormand also uses the seventh stage to compute the
eighth and ninth stages, whereas Owren only uses the first six stages.
Still another method, cited by Montenbruck and Gill [25], is the Keplerian
interpolation. This method uses the osculating Keplerian orbit which matches the
state vector at the beginning of the integration step. In comparison with polynomial
interpolations discussed previously, this method requires greater computational
effort, but is not limited to a single integration step. A full account of various
interpolants can be found in [25].
The expressions shown above may be generalised as follows. The solution of
order p and the solution of order (p + 1) for an embedded Runge–Kutta formula
with s stages are, respectively, (see, e.g., [17, 26]):
P
s P
s
zn þ 1 ½p ¼ zn þ h b i ki zn þ 1 ½p þ 1 ¼ zn þ h bi ki
i¼1 i¼1
where
!
iP
1
ki ¼ f tn þ ci h; zn þ h aij kj ði ¼ 1; 2; . . .; sÞ
j¼1
Following Enright [27], the continuous methods use m additional stages ks+1,
ks+2, …, ks+m and polynomials bi(r) of degree less than or equal to p, for i = 1, 2,
…, (s + m), such that the interpolating approximation can be expressed by
z½p þ 1 ðsÞ ¼ z½p þ 1 ðtn þ rhÞ ¼ zn þ h½b1 ðrÞk1 þ b2 ðrÞk2 þ þ bs þ m ðrÞks þ m
where, as shown above, the following relations hold:
X
i1
ci ¼ aij
j¼1
r 2 ½0; 1
s tn
r¼
h
bi ð1Þ ¼ bi for i ¼ 1; 2; . . .; s
bi ð1Þ ¼ 0 for i ¼ s þ 1; s þ 2; . . .; s þ m
In other words, a continuous embedded Runge–Kutta method RKp(p + 1)s may

be represented by means of the following Butcher tableau, where the continuous
weights shown in the last row (b1, b2, …, bs+m) are polynomials in r, with
s tn
r¼
h
There are various methods to construct continuous extensions to discrete Runge–

Kutta formulae. One of them is based on the classical theory of interpolation.
Let us consider a generic discrete Runge–Kutta formula of order q. Following
Dormand [24], the continuous extension of the same local order is provided by a
Hermite interpolant p(s), of degree greater than or equal to q, based on the end
points of the discrete integration step and an appropriate number of internal points
or derivatives having an appropriate order of accuracy. For example, if we take
q = 5, the interpolant must respect the following conditions
pð t n Þ ¼ z n p0 ðtn Þ ¼ z0 ðtn Þ
pð t n þ 1 Þ ¼ z n þ 1 p0 ðtn þ 1 Þ ¼ f ðtn þ 1 ; un þ 1 Þ
p0 ðtn þ ri Þ ¼ f ðtn þ ri ; un þ ri Þ i ¼ 1; 2
where un+ri is an estimate of z(tn+ri), obtained from a fourth-order interpolant,

which is presumed to be available. These six conditions determine the six coeffi-
cients of a polynomial of the fifth degree.
Let p(s) = a0r0 + a1r1 + a2r2 + + a5r5 be the expression of the fifth-degree
Hermite polynomial, where r = (s − tn)/h. The coefficients a0, a1, a2, …, a5 of this
polynomial can be determined by solving a system of linear algebraic equations. If

a cubic Hermite polynomial is used, there is no need for internal points. An
application of this concept has been made by Montenbruck and Gill [25] to
on-board navigation systems. The state vector is approximated by the following
cubic Hermite polynomial
zðtn þ rhÞ ¼ a0 ðrÞzðtn Þ þ a1 ðrÞhz0 ðtn Þ þ a2 ðrÞzðtn þ 1 Þ þ a3 ðrÞhz0 ðtn þ 1 Þ
where
a0 ðrÞ ¼ ðr 1Þ2 ð2r þ 1Þ

a1 ðrÞ ¼ rðr 1Þ2
a2 ðrÞ ¼ r2 ð3 2rÞ
a3 ðrÞ ¼ r2 ðr 1Þ
This polynomial matches the state vectors z(tn) and z(tn+1) and the respective
derivatives z′(tn) and z′(tn+1) at the beginning and at the end of the integration step
h. This requires the knowledge of the derivative z′(tn+1) = f[tn + h, z(tn+1)], which is
available from the discrete formula. Using this information and the position (x),
velocity (x′) and acceleration (x″) vectors at the beginning (tn) and at the end (tn+1)
of the integration step, we have
xðtn þ rhÞ ¼ a0 ðrÞxðtn Þ þ a1 ðrÞhx0 ðtn Þ þ a2 ðrÞh2 x00 ðtn Þ þ a3 ðrÞxðtn þ 1 Þ

þ a4 ðrÞhx0 ðtn þ 1 Þ þ a5 ðrÞh2 x00 ðtn þ 1 Þ
where
a0 ðrÞ ¼ 110r3 þ 15r4 6r5

a1 ðrÞ ¼ r6r3 þ 8r4 3r5
a2 ðrÞ ¼ 1=2ðr2 3r3 þ 3r4 r5 Þ
a3 ðrÞ ¼ 1 a0 ðrÞ ¼ 10r3 15r4 þ 6r5
a4 ðrÞ ¼ 4r3 þ 7r4 3r5
a5 ðrÞ ¼ 1=2ðr3 2r4 þ r5 Þ
By differentiating the expression of x(tn + rh), we have
x0 ðtn þ rhÞ ¼ a0 ðrÞxðtn Þ=h þ a1 ðrÞx0 ðtn Þ þ a2 ðrÞhx00 ðtn Þ þ a3 ðrÞxðtn þ 1 Þ=h
þ a4 ðrÞx0 ðtn þ 1 Þ þ a5 ðrÞhx00 ðtn þ 1 Þ
where
a0 ðrÞ ¼ 30r2 þ 60r3 30r4

a1 ðrÞ ¼ 118r2 þ 32r3 15r4
a2 ðrÞ ¼ 1=2ð2r9r2 þ 12r3 5r4 Þ
a3 ðrÞ ¼ a0 ðrÞ ¼ 30r2 60r3 þ 30r4
a4 ðrÞ ¼ 12r2 þ 28r3 15r4
a5 ðrÞ ¼ 1=2ð3r2 8r3 þ 5r4 Þ
By combining the data from two consecutive integration steps (one going from tn
to tn+1 and the other going from tn+1 to tn+2), we have a fifth-degree Hermite
polynomial approximating the state vector

þ a4 ðrÞzðtn þ 2 Þ þ a5 ðrÞhz0 ðtn þ 2 Þ
where
a0 ðrÞ ¼ 1=4ðr1Þ2 ðr2Þ2 ð1 þ 3rÞ

a1 ðrÞ ¼ 1=4rðr1Þ2 ðr2Þ2
a2 ðrÞ ¼ r2 ðr2Þ2
a3 ðrÞ ¼ ðr1Þr2 ðr 2Þ2
a4 ðrÞ ¼ 1=4r2 ðr1Þ2 ð73rÞ
a5 ðrÞ ¼ 1=4ðr2Þr2 ðr1Þ2
Otherwise, the method of divided differences may be used to determine the

coefficients of the polynomial. Enright et al. [26, 28] have proposed a method,
called bootstrapping process, which constructs continuous extensions by means of
interpolating polynomials of increasing degree, as will be shown below.
The bootstrapping process extends the concept of Hermite interpolation by
means of additional function values comprised within the integration steps; these
values are used to compute approximations of higher order. Montenbruck and Gill
[25] cite two examples to illustrate this concept. As a first example, previously
given by Enright et al. [28], they consider the discrete Runge–Kutta–Fehlberg RK4
(5)6 formula pair (shown here in Sect. 6.3). A fifth-order approximation zn+0.6 is
computed at the internal point tn + 0.6 h. This approximation is expressed by [25,
28]:

1559 153856 68107 243 2106
zn þ 0:6 ¼ zn þ h k1 þ k3 þ k4 k5 k6
12500 296875 2612500 31250 34376
and is introduced in a fourth-degree Hermite polynomial

þ a4 ðrÞzn þ 0:6
where

2 5 11
a0 ðrÞ ¼ ðr 1Þ 1 r rþ1
3 3

5
a1 ðrÞ ¼ rðr 1Þ2 1 r
3

3 5
a2 ðrÞ ¼ r2 r ð9r 11Þ
4 4

5 3
a3 ðrÞ ¼ r ðr 1Þ r
2
2 2
ðr 1Þ2
a4 ðrÞ ¼ 625 r2
36
As a second example, Montenbruck and Gill cite the application made by

Shampine to the fifth-order Runge–Kutta method of Dormand and Prince.
In this case, a fifth-order result is computed at the internal point tn + 1=2h. This
approximation is expressed by [25]:

1 7157 709256 10825 220887
zn þ 0:5 ¼ zn þ h k1 þ k3 þ k4 k5
2 37888 82362 56832 2008064

80069 107 5
þ k6 k7 k8
1765344 2627 37
where k8 is expressed by

1 33728713 30167461 7739027
k7 ¼ f tn þ h; zn þ h k1 þ 2k2 k3 þ k4
2 104693760 21674880 17448960

19162737 26949
k5 k7
123305984 363520
This approximation and the endpoints zn and zn+1 are used to construct a
fourth-degree polynomial over the integration interval.
6.10 Symplectic Explicit Special Nyström Methods
There are time-dependent processes or physical phenomena which are

non-dissipative; that is, there is no loss of energy as the process or phenomenon
evolves with time. Such cases as these, arising for example in astrophysics or
molecular dynamics, can be mathematically modelled by means of a Hamiltonian
system of differential equations of motion [29]:
p0 ¼ f ðp; q; tÞ
q0 ¼ gðp; q; tÞ
where the vector-valued functions f {f1 f2 … fn}T and g {g1 g2 … gn}T are
such that
@H @H @H
f1 ¼ f2 ¼ fn ¼
@q1 @q2 @qn
@H @H @H
g1 ¼ g2 ¼ gn ¼
@p1 @p2 @pn
for sufficiently smooth real functions H = H(p, q, t), where q {q1 q2 … qn}T is
the n 1 column vector of the generalised co-ordinates and p {p1 p2 … pn}T is
the n 1 column vector of the generalised momenta. When the Hamiltonian
H does not depend explicitly on time, we have an autonomous Hamiltonian system.
Numerical methods of integration for solving these systems of equations should
likewise be non-dissipative. Such traditional methods as those of the Runge–Kutta
type are not the most appropriate methods for the purpose of integrating
Hamiltonian systems, because these methods introduce numerical dissipative per-
turbations, against which Hamiltonian systems are not structurally stable. In other
words, a Hamiltonian system, integrated by means of a traditional numerical
method, becomes at length a dissipative (i.e., non-Hamiltonian) system, which
behaves in the long term in a manner which differs from that which is proper of its
nature. This necessity has led some authors to study symplectic integration methods
for Hamiltonian systems. These methods preserve the Hamiltonian structure of the
system by arranging that each step of the integration should be a canonical or
symplectic transformation.
The equations of motion written above may also be written in compact form by
defining a 2n 1 column vector x such that x [q p]T. By using this notation, the
equations of motion written above may be written as follows
@H
x0 ¼ J
@x
where J is the following 2n 2n skew-symmetric matrix

6.10 Symplectic Explicit Special Nyström Methods 935

0 I
J¼
I 0
and I and 0 are, respectively, the n n identity and the zero matrices. The pre-
ceding expression is the so-called symplectic notation for Hamilton’s equations.
Following Chou and Sharp [30], single-step numerical methods are called
symplectic if they preserve the structure of the space of variables (q, p), thus
reproducing the main qualitative property of the solution. This means preservation
of volume in the phase space (which represents all possible values of the position
and momentum variables), that is, the Liouville property.
Following Okunbor and Lu [31], let M be the Jacobian of map for the integration
step h, that is, for the transformation from q(nh) and p(nh) to, respectively,
q[(n + 1)h] and p[(n + 1)h]. In order for a numerical integration method to be
symplectic, the method must not only preserve some qualitative features of the
dynamical system represented, but also be consistent with the latter. Consequently,
Okunbor and Lu define a method as symplectic for any h and any Hamiltonian
system for which it is applicable, if its Jacobian matrix M satisfies the condition
MTJM = J.
Numerical experiments have shown that symplectic methods are more efficient
than non-symplectic methods for long intervals of integration. Symplectic Runge–
Kutta methods, for such general cases as those indicated above, are necessarily
implicit, because the matrix A {aij} of the coefficients aij in the Butcher tableau
is a not a strictly lower triangular matrix. However, in case of separable
Hamiltonians of the form
H ðp; qÞ ¼ T ðpÞ þ V ðqÞ
where T(p) and V(q) are called in mechanics, respectively, kinetic energy and
potential energy, explicit Runge–Kutta–Nyström methods are called symplectic if
the coefficients of the method respect certain conditions, which will be shown
below. Such is the case with most physical problems, where the vector field is
separable in two exactly solvable parts. In these problems, the exact solution can be
approximated by a composition of flows associated with each part of the vector field
[32]. For example, in the case of a body of unit mass orbiting around the Earth (i.e.,
the plane two-body problem which will be considered in Sect. 6.11), the
Hamiltonian is the total (kinetic plus potential) energy of the body
1 l 1 2 l
H ¼ v2 ¼ px þ p2y 12
2 r 2
q2x þ q2y
where r and v are, respectively, the radius vector and the velocity of the body, and l
is the gravitational parameter of the Earth (l is equal to G times the mass of the
Earth, where G is the gravitational constant).
When a Hamiltonian H(p, q) is separable, the autonomous Hamiltonian system
of differential equations p′ = f(p, q) and q′ = g(p, q) can be written in the form
x00 ¼ f ðxÞ
where the function f(x) is the gradient of −V = −V(x), x q is the vector of the
generalised co-ordinates, and x′ p is the vector of the generalised momenta [33].
In this case, explicit special Nyström methods (described in Sect. 6.8) can be used
to generate approximations of the following type
X
s
xn þ 1 ¼ xn þ hx0 n þ h2 b^ i ki
i¼1
X
s
x0 n þ 1 ¼ x0 n þ h bi ki
i¼1
where
!
0
P
s
ki ¼ f tn þ ci h; xn þ ci h x n þ h 2
aij kj ði ¼ 1; 2; . . .; sÞ
j¼1
and where b^1, b^2, …, b^s and b1, b2, …, bs denote, respectively, the weights to be
used in the formula for xn+1 and the weights to be used in the formula for x′n+1. An
explicit special Nyström method is called symplectic if the following two condi-
tions are satisfied
b^ i ¼ ð1 ci Þbi ði ¼ 1; 2; . . .; sÞ
aij ¼ ðci cj Þbj ðj ¼ 1; 2; . . .; i 1; i ¼ 2; 3; . . .; sÞ
Consequently, if the coefficients ci and bî are known, then the remaining
coefficients of an explicit special Nyström method are also known. As has been
shown in Sect. 6.8, in order for a method to be FSAL (i.e., such that the First stage
is the Same As the Last), the coefficients of the method must satisfy the conditions
c1 = 0, cs = 1, and asi = bî (i = 1, 2, …, s − 1).
Hence, if c1 = 0, cs = 1, and the method is symplectic, then the condition
asi = bî (i = 1, 2, …, s − 1) is automatically satisfied.
Chou and Sharp have found a seven-stage, fifth-order, FSAL, explicit, special
Nyström method which has the following coefficients (from [30]):
c1 ¼ 0:0
c2 ¼ 0:2179621390175646
c3 ¼ 0:4424703708255242
c4 ¼ 0:1478460559438898 101
c5 ¼ 0:3400000000000000
c6 ¼ 0:7000000000000000
c7 ¼ 1:0
b1 ¼ 0:6281213570268329 101
b2 ¼ 0:3788983131252575
b3 ¼ 0:2754528515261340
b4 ¼ 0:1585299574780513 102
b5 ¼ 0:1785704038527618
b6 ¼ 0:3479995834198831
b7 ¼ 0:1149928196535844
The remaining coefficients bî and aij result from the conditions indicated above,
that is,
b^ i ¼ ð1 ci Þbi ði ¼ 1; 2; . . .; 7Þ
aij ¼ ðci cj Þbj ðj ¼ 1; 2; . . .; i 1; i ¼ 2; 3; . . .; 7Þ
As reported by Chou and Sharp, the Euclidean norms of the six error coefficients
for the solution formula (xn+1) and the derivative formula (x′n+1) are, respectively,
4.0 10−4 and 4.1 10−4. The results of the performance evaluations performed
by Chou and Sharp will be given in Sect. 6.11.
Other symplectic explicit special Nyström methods are due to:
(a) Calvo and Sanz-Serna (five-stage, fourth-order, FSAL), the coefficients of
which are given below (from [33])
c1 ¼ 0:0
c2 ¼ 0:2051776615422863869
c3 ¼ 0:6081989431465009739
c4 ¼ 0:4872780668075869657
c5 ¼ 1:0
b1 ¼ 0:0617588581356263250
b2 ¼ 0:3389780265536433551
b3 ¼ 0:6147913071755775662
b4 ¼ 0:1405480146593733802
b5 ¼ 0:1250198227945261338
(b) Okunbor and Skeel (five-stage, fifth-order, non-FSAL), the coefficients of

c1 ¼ 0:69883375727544694289
c2 ¼ 0:20413810365459889029
c3 ¼ 1:02055757000418534370
c4 ¼ 0:36292800323075291580
c5 ¼ 0:30508610893167564804
b1 ¼ 0:40090379269664777606
b2 ¼ 0:95997088013412390506
b3 ¼ 0:08849515812721633901
b4 ¼ 1:22143909234910252870
b5 ¼ 1:67080892330709041000
(c) Okunbor and Skeel (seven-stage, sixth-order, non-FSAL), the coefficients of

c1 ¼ 1 c7
c2 ¼ 1 c6
c3 ¼ 1 c5
c4 ¼ 0:5
c5 ¼ 1:43531315933193655010
c6 ¼ 0:24517048359575719767
c7 ¼ 0:88961673353684493504
b1 ¼ b7
b2 ¼ b6
b3 ¼ b5
b4 ¼ 0:00024286040977501724
b5 ¼ 0:08191385007043372004
b6 ¼ 0:23158642248235284281
b7 ¼ 0:64955114220703161414
(d) Calvo and Sanz-Serna (thirteen-stage, eighth-order, FSAL), the coefficients of

which are given below (from [33, 35])
c1 ¼ 0:0
c2 ¼ 0:60715821186110352503
c3 ¼ 0:96907291059136392378
c4 ¼ 0:10958316365513620399
c5 ¼ 0:05604981994113413605
c6 ¼ 1:30886529918631234010
c7 ¼ 0:11642101198009154794
c8 ¼ 0:29931245499473964831
c9 ¼ 0:16586962790248628655
c10 ¼ 1:22007054181677755238
c11 ¼ 0:20549254689579093228
c12 ¼ 0:86890893813102759275
c13 ¼ 1:0
b1 ¼ 1=2c2
bi ¼ 1=2ðci þ 1 ci1 Þ i ¼ 2; 3; . . .; 12
b13 ¼ 1=2ð1 c12 Þ
Sharp and Vaillancourt [35] have made performance evaluations for all the
methods shown in the present paragraph.
6.11 Performance Comparison for

Runge–Kutta(–Nyström) Methods
In order to evaluate the performance offered by a Runge–Kutta method in problems

of orbit computation, some test cases have been suggested by various authors. From
the point of view of astrodynamics, it is particularly important to mention the plane
two-body problem, proposed by Hull et al. in 1972 (see [17]), which concerns the
motion of a satellite around the Earth.
The equations of motion, in terms of Cartesian co-ordinates x and y, are
x y
x00 ¼ y00 ¼
r3 r3
where r = (x2 + y2)½ is the radius vector, and
x0 ¼ 1 e y0 ¼ 0
12
0 0 1þe
x0 ¼0 y0 ¼
1e
are the initial conditions (t = 0). As is well known, the analytical solution is an
elliptical orbit, with semi-major axis a = 1 and eccentricity e, starting at perigee.
For the purpose of integrating numerically the equations of motion, they are
rewritten below in terms of the state vector and its time derivatives.
Let z [z1 z2 z3 z4]T [x y x′ y′]T be the 4 1 state vector, whose components
are the satellite position and velocity. The four scalar equations of motion repre-
senting this problem are then
z1 z2
z0 1 ¼ z3 z0 2 ¼ z4 z0 3 ¼ 32 z0 4 ¼ 32
z21 þ z22 z21 þ z22
with the corresponding initial values

12
1þe
z10 ¼ 1 e z20 ¼ 0 z30 ¼ 0 z40 ¼
1e
Hull et al. have defined five test problems with eccentricity values ranging from
0.1 to 0.9 in steps of 0.2. The final time was tF = 20 in all cases, corresponding to
20/(2p) 3.2 revolutions around the Earth. The analytical solution is
1
xðtÞ z1 ðtÞ ¼ cos Æ e yðtÞ z2 ðtÞ ¼ ð1 e2 Þ2 sin Æ
1
sin Æ ð1 e2 Þ2 cos Æ
x0 ð t Þ z 3 ð t Þ ¼ y0 ðtÞ z4 ðtÞ ¼
1 e cos Æ 1 e cos Æ
where Æ is the eccentric anomaly, which is computed by solving iteratively

Kepler’s equation t = Æ − e sin Æ.
By comparing the analytical solution with each of those computed numerically,
it is possible to compute the global error committed in the numerical integration.
The results of the comparison can be found in Montenbruck’s papers [17, 36].
A summary of these results is given below.
The methods considered include Runge–Kutta methods for (systems of)
first-order differential equations and Runge–Kutta–Nyström methods for
second-order differential equations with or without velocity-dependent terms. The
particular methods are:
• RK5(4)7FM by Dormand and Prince (see [15] and also here, Sects. 6.7 and 6.9),
which is a seven-stage method of the fifth order with an embedded formula of
the fourth order for error control;
• RK8(7)13M by Prince and Dormand (see [16] and also here, Sect. 6.7), which is
a thirteen-stage method of the eighth order with an embedded formula of the
seventh order for error control;
• RK7(8)13 by Fehlberg (see [6] and also here, Sect. 6.4), which is similar to RK8
(7)13M, but here the seventh-order formula is used for integration and the
eighth-order formula is used for error control;
6.11 Performance Comparison for Runge–Kutta(–Nyström) Methods 941
• RKN7(6)9 by Dormand and Prince (see [18]), an embedded nine-stage,

seventh-order method, which unlike the methods indicated above integrates
directly the special, second-order equation x″ = f(t, x);
• RKN(x′)7(8) by Fehlberg (see [37]), a thirteen-stage Runge–Kutta–Nyström
method, which integrates directly the general, second-order equation x″ = f(t, x, y);
• RKN11(10)17 by Filippi and Gräf (see [19 and also here, Sect. 6.8), which is a
seventeen-stage, eleventh-order special Nyström method; and
• RKN12(10)17M by Dormand and Prince (see [24]), which is a seventeen-stage,
12(10)th-order method.
The results found by Montenbruck show that high-order methods are necessary
to reach accuracies equal to or greater than 10−8. A comparison between RK7(8)13
(Fehlberg) and RK8(7)13M (Prince and Dormand) shows the advantage of the latter
method, which uses the eighth-order formula to compute the solution and the
seventh-order formula to control the step size. This method has also been found
advantageous in comparison with Fehlberg’s RKN(x′)7(8).
Higher-order codes, such as RKN11(10)17 by Filippi-Gräf and RKN12(10)17M
by Dormand-Prince, have a greater efficiency than those cited above, even when the
accuracy required by the user is low.
Further performance comparisons than these have been made by Dormand and
Prince [18], using the same plane two-body problem proposed by Hull et al. (cited
above). Dormand and Prince have used an embedded fourth-order Runge–Kutta
method (RK4(3)T) with k = 1/10 and an embedded fifth-order Runge–Kutta
method (RK5(4)T) with k = 1/60; they have also used the embedded seventh-order
RKN7(6)T special Nyström method with k = 1/20.
The coefficients for RKN7(6)T have been given in the corresponding table of
Sect. 6.8. Those for RK4(3)T and RK5(4)T are given below (from [18]).
Coefficients for RK4(3)T Runge–Kutta method

Source Dormand and Prince [18]
Coefficients for RK5(4)T Runge–Kutta method are

Source Dormand and Prince [18]
–
–
– –
Dormand and Prince have also considered the embedded Runge–Kutta–Fehlberg

RK4(5)6, RK5(6)8, and RK7(8)13 methods, which have been cited several times
here; in addition to these, they have considered the special Nyström method due to
Bettis ([11] and also here, Sect. 6.6). The results are shown by Dormand and Prince
[18] in terms of efficiency curves (maximum global error versus function evalua-
tions), separately for Runge–Kutta and Runge–Kutta–Nyström methods. Such
results are summarised below.
Dormand and Prince find that RK7(8)13, used in eighth-order mode, works
better than other methods for global errors less than 10−3 approximately. Otherwise,
they recommend to use either RK4(5)6 in fifth-order mode or RK5(4)T.
The special Nyström methods, when applicable, are to be preferred to the cor-
responding Runge–Kutta methods of the same order. The method of Bettis has been
found by Dormand and Prince the most efficient for low values of accuracy,
whereas for high accuracy, the RKN7(6)T has been found to be the best formula.
The performance evaluations made by Chou and Sharp [30] for symplectic
explicit special Nyström methods, described in Sect. 6.10, are also based on the
plane two-body problem, proposed by Hull et al. in 1972 (see [17]). In terms of
p and q variables, the Hamiltonian is
1 2 1
H¼ px þ p2y 12
2
q2x þ q2y
6.11 Performance Comparison for Runge–Kutta(–Nyström) Methods 943
and the initial conditions are

12
1þe
qx0 ¼ 1 e qy0 ¼ 0 px0 ¼ 0 py0 ¼
1e
The results of [30] are expressed by diagrams plotting the number of evaluations
of the integrand function against the Euclidean norm of the end-point global error.
The eccentricity values considered by Chou and Sharp are three (0.3, 0.5, and 0.7)
over an integration interval of 10,000 with step sizes of 2−2, 2−3, …, 2−7. The
fifth-order method proposed by Chou and Sharp (the coefficients of which method
are given here in Sect. 6.10) is found advantageous by them in comparison with
other fifth-order methods cited in [20], except at small global errors for a value of
eccentricity equal to 0.3.
Likewise, Blanes and Moan [32] have made performance evaluations based on
the same two-body problem for fourth- and sixth-order explicit symplectic methods
and values of eccentricity ranging from 0.1 to 0.8. They have integrated the dif-
ferential equations over ten orbital periods and have measured the average error in
position, along the last two periods, for different values of the eccentricity. The
results are given in graphs plotting the error (in logarithmic scale) against eccen-
tricity, with each curve corresponding to a single method. The methods tested by
them are symplectic methods with order ranging from four to six.
Sharp and Vaillancourt [35] have made performance evaluations for all the
explicit symplectic methods shown in Sect. 6.10, by means of long n-body simu-
lations of the Solar System. The step sizes used by them are equal to 4, 8, 16, 32,
and 64 days, over a total integration time of one million years for the Sun and the
gas giants (Jupiter, Saturn, Uranus, and Neptune), 22,000 years for the Sun and the
nine planets, and 11,000 years for three satellites of Saturn. The diagrams plot, in a
log-log scale, the Euclidean norm of total error against time, for each of the step
sizes indicated above. A diagram is shown for each of the symplectic methods
considered by the authors. The results found by them are summarised below. In the
case of the eighth-order method due to Calvo and Sanz-Serna [33] applied to the
planetary system comprising the Sun and the gas giants, there seems to be an
approximately linear growth of the total error norm with time, for step sizes of 32
and 64 days, whereas the total error norm grows faster than linearly for step sizes of
4, 8, and 16 days. The comparative evaluation, which took in consideration all of
the five symplectic methods described in Sect. 6.10, shows that:
• for large step sizes, the total error grows as time raised to a power of 1 + e,
where ||e|| 1, except for the sixth-order method of Okunbor and Skeel in the
simulation of the Sun and the nine planets, in which case the exponent is
approximately 1.18;
• for small step sizes, the total error grows as time raised to an exponent greater
than unity;
• the solution for the higher-order method of Calvo and Sanz-Serna with a step
size of 2h has been found at least as accurate as that for the sixth-order method
of Okunbor and Skeel with a step size of h, with consequent advantages of the
former over the latter in terms of both number of function evaluations and (to a
lesser extent) CPU time;
• the fifth-order method of Chou and Sharp has been found more efficient than the
eighth-order method of Calvo and Sanz-Serna for large step sizes and less
efficient for small step sizes;
• the fourth-order method of Calvo and Sanz-Serna has been found less efficient
than their eighth-order method for all the step sizes used in the simulations.
Calvo and Sanz-Serna [33] have made performance evaluations on the same
plane two-body problem as that shown above. They have considered their sym-
plectic methods of, respectively, fourth order (Sect. 6.10, subparagraph a) and
eighth order (Sect. 6.10, subparagraph d), another symplectic eighth-order method
due to Yoshida [38], and a non-symplectic method due to Dormand et alii [39].
Their experiments are limited to a value of eccentricity equal to 0.5. The results are
given in two diagrams, which refer to an integration interval of, respectively,
810 2p and of 21870 2p, where the total error is plotted against the number of
function evaluations. Both of these diagrams indicate a greater efficiency of the
symplectic high-order formula of Calvo and Sanz-Serna in comparison with the
correspondent formula due to Yoshida. The latter method has been found more
efficient than the fourth-order symplectic method of Calvo and Sanz-Serna only
when small errors are required; otherwise, the opposite has turned out to be true.
The non-symplectic formula due to Dormand et alii has been found more efficient
than any of the symplectic formulae tested.
6.12 Bulirsch-Stoer Methods
These methods are characterised by automatic order and step-size control. They
estimate the error by: (a) computing a solution over an interval by means of the
same method used repeatedly with a varying number of steps; and (b) extrapolating
the polynomial which fits through the computed solutions. Such methods are based
on Gragg’s [40] modified midpoint rule (which has only even powers of the step
size in its error expansion) and Richardson’s deferred approach to the limit [41, 42].
The modified midpoint rule is based on the following formulae.
6.12 Bulirsch-Stoer Methods 945
Let x′= f(t, x) be a differential equation defined in an interval a t b with

the initial condition x0 = x(t0). Let uN be an approximate solution of the true
solution x at t = tN, that is, uN x(tN). We want to advance from uN to uN+1 at
tN + H, by taking a set of n sub-steps each of size h, that is,
H
h¼
n
Following Bulirsch and Stoer [43], the modified midpoint rule is expressed by
the operations indicated below:
z 0 ¼ uN
z1 ¼ uN þ hf ðtN ; uN Þ
zm þ 1 ¼ zm1 þ 2hf ðt þ mh; zm Þ ðm ¼ 1; 2; . . .; n 1Þ
uN þ 1 ¼ 1=2½zn þ zn 1 þ hf ðtN þ H; zn Þ
where z0, z1, … are intermediate approximations which proceed in steps of h, and
uN+1 is the final approximation to x(tN + H).
The computed result is extrapolated to the value that would have been obtained
if the step size had been much smaller than it really was. In particular, we want to
extrapolate the computed result to zero step size. The principle of this method is
that, when the step size is decreased, the computed result becomes more accurate
than the result computed with the previous step size. Consequently, we take a large
step H, comprising many sub-steps taken by means of the modified midpoint rule.
Each sub-step produces an intermediate estimate (z0, z1, …, zn) of the solution. The
results computed by means of these sub-steps converge towards the result which
would have been obtained with a step size equal to zero. Therefore, the converging
results are used in a polynomial approximation of the results as a function of the
step size, to obtain the result relating to zero step size.
The deferred approach to the limit is a method to improve an approximation that
depends on a step size. In general terms, this method may be presented (see [44]) as
follows. We want to compute, in a finite process, a certain quantity, for example, an
integral I, by means of approximations using different step sizes. Let I(h) be an
approximation of I that depends on a step size h with an error formula of the form
I I ð h Þ ¼ a1 h 2 þ a 2 h 4 þ a 3 h 6 þ
where a1, a2, a3, … are unknown constants. The exact value sought can be
expressed as

I ¼ I ð hÞ þ a1 h2 þ a2 h 4 þ O h6
Using step sizes h and h/s for some s, we can write

I ¼ I ð hÞ þ a1 h2 þ a2 h 4 þ O h6
2 4
h h h
I¼I þ a1 þ a2 þ O h6
s s s
By multiplying all terms of the second equation by s2 and subtracting the cor-
responding terms of the first equation, we have

h
s 1 I ¼s I
2
I ð hÞ þ O h6
2
s
which can be solved for I to yield

h
s2 I I ð hÞ
s
I¼ þ O h6
s 1
2
By this process, we have achieved a better approximation of I by subtracting the

largest term in the error which was O(h6). The process may be repeated to remove
more error terms to get better approximations. Lynch [45] has applied this method
to the problem of evaluating p by inscribing regular polygons having an increasing
number of sides in a unit circle. A further application of Richardson extrapolation is
Romberg integration (see, e.g., [46]), which applies Richardson extrapolation to the
trapezium rule.
As shown above, the Bulirsch-Stoer method is based on: (1) Richardson’s idea
of expressing the result as a function of the step size and using that function to
extrapolate to the result corresponding to zero step size; (2) Gragg’s idea of
expressing the error term, by means of proper manipulation, by only even powers of
the step size, so that at each stage of the polynomial approximation, it is possible to
cancel two powers at a time of the error term; and (3) Bulirsch-Stoer’s idea of using
a rational approximation, instead of Gragg’s polynomial approximation, in order to
eliminate the difficulties of the Taylor series in the vicinity of poles and singularities
in the complex plane. Another improvement due to Bulirsch and Stoer is their
method of reducing the step size by means of a sequence suggested by them, as will
be shown below.
According to the Bulirsch-Stoer method, we compute an approximate solution
uN+1 at tN + H by means of the modified midpoint rule, for various values of n. When
we have a few of such values, we fit the solution to an analytic form and evaluate it at
h = 0, in accordance with the method of the deferred approach to the limit.
The sequence of individual attempts, performed according to the modified
midpoint rule, is made by using increasing values of n. Common choices for the
extrapolation sequence are shown below.
(a) Bulirsch sequence [43]:
n ¼ 2; 4; 6; 8; 12; 16; 24; 32; . . . nj ¼ 2nj2 ðj

4Þ
(b) Deuflhard sequence [47]:
n ¼ 2; 4; 6; 8; 10; 12; 14; . . . nj ¼ 2j
(c) harmonic sequence:
n ¼ 1; 2; 3; 4; 5; 6; 7; 8; . . . nj ¼ j
According to Gragg [40], the difference between the approximate solution

uN+1 = uN+1(h) and the exact solution x(tN + H) can be expressed by an asymptotic
expression in h2
uN þ 1 ðhÞ ¼ xðtN þ H Þ ¼ e2 h2 þ e4 h4 þ e6 h6 þ
for sufficient smooth functions and even values of n. The error coefficients ej depend
on tN and H, but are independent of h.
The formulae used for the sub-steps have a low degree of approximation.
However, the degree of approximation is improved by repeating the integration
with a different step size h, which leads to
uN þ 1 ðhÞ ¼ xðtN þ H Þ ¼ 2 h2 þ 4

h4 þ 6
h6 þ
Now, by means of the expression
h2 uN þ 1 ðhÞ h2 uN þ 1 ðhÞ 2 2

u N þ 1 ¼ ¼ x ð tN þ H Þ þ 4 O h
h
h h2
2
we can eliminate the leading term (e2h2) of the error. This results in a reduction of
the error by two orders. The extrapolation uses the Aitken-Neville technique, which
is described below. The first column
T11 ¼ uN þ 1 ð h1 Þ
T21 ¼ uN þ 1 ð h2 Þ
T31 ¼ uN þ 1 ð h3 Þ
T41 ¼ uN þ 1 ð h4 Þ
..
.
of the triangular array shown below results from using Gragg’s modified midpoint
rule for increasing sub-steps:
T11
&
T21 ! T22
& &
T31 ! T32 ! T33
& & &
T41 ! T42 ! T43 ! T44
..
.
If we use a polynomial extrapolation, each entry in the table (or array) shown
above is a linear combination of the entries to the left and upper left of it, as
indicated by the arrows.
Consequently, the following (second, third, etc.) columns of the table are filled
by using the recurrence relation

h hjk þ 1 Tj;k1 ðhÞ h hj Tj1;k1 ðhÞ
Tj;k ðhÞ
hj hjk þ 1
The value corresponding to extrapolation at h = 0 is
hjk þ 1 Tj;k1 ðhÞ þ hj Tj1;k1 ðhÞ Tj;k1 Tj1;k1

Tj;k ð0Þ ¼ ¼ Tj;k1 þ 2
hj hjk þ 1 nj
njk þ 1 1
After each successive n is tried, a polynomial interpolation is attempted.

Adjacent diagonal entries of the Aitken-Neville array are used for error estimates. If
the error is greater than the tolerance, the next n in the sequence is tried.
Error control is performed by using the procedure suggested by Deuflhard [47];
that is, we estimate the error by means of the sub-diagonal entries (Tk+1,k) of the
extrapolation array

ek þ 1;k ¼
Tk þ 1;k Tk þ 1;k þ 1
and impose the following condition
ek þ 1;k \e
where ek+1,k is the local truncation error of column k, and e is the tolerance specified
by the user.
Such estimates can be shown [48] to be of the order of O(H2k−1). For an even
sequence of n, the order of the method is 2 k − 1, and the new step size may be
estimated as suggested by Deuflhard [47]:
2k1
1

Hk ¼ H
k þ 1;k
Montenbruck and Gill [17] suggest to include a safety factor equal to 0.9 in the
computation of the new step size, that is,
2k1
1

Hk ¼ 0:9H
k þ 1;k
6.13 Multi-step Methods
Given a differential equation in vector form x′= f(t, x) and an associated initial
condition x0 = x(t0), the single-step methods shown in the preceding paragraphs
compute the dependent-variable vector xi+1 at the next mesh point ti+1 by means of
information coming only from xi and the integrand function f(t, x) evaluated in
some points of the subinterval [ti, ti+1] which are placed at certain fractions of the
current step-size hi = ti+1 − ti of the independent variable t. Such methods are,
therefore, local methods.
On the contrary, the multi-step methods compute xi+1 by means of information
from several previous mesh points of the solution. A category of multi-step
methods is based on quadrature formulae. These methods replace the integrand
function by an interpolating polynomial, then predict the next value of the solution
by means of an explicit formula, and then correct the predicted value by means of
an implicit formula. Such is the case with the Adams, Milne-Simpson, and Nyström
formulae. Multi-step methods of another category, which are called backward
differentiation formulae and are not considered in the present book, perform iter-
ations which use implicit formulae to compute the next value xi+1 of the solution.
The multi-step methods of both categories are, therefore, global methods.
Let the values of the dependent variable x be known in k points of the interval of
integration a t b. The multi-step methods compute the next value xi+k as a
function of the preceding values xi+k−1, xi+k−2, …, xi (and, of course, of the inte-
grand function). Such methods are also called k-step methods. In particular, if
k were equal to one, the corresponding method would be a single-step method, as
those considered in the preceding paragraphs. A formula is called explicit, if xi+k is
expressed as a function of the preceding values only, or implicit, if xi+k is expressed
as a function of the preceding values and of xi+k itself (appearing both on the
left-hand and on the right-hand sides of the equation). The integration interval [a, b]
is assumed as divided into segments of equal length, for the purpose of avoiding the
difficulties (which will be considered in the next paragraphs) connected with a

change of step size. It is also assumed that there is a linear relation between xi+k and
fi+k (for any value of k), so that a multi-step method can be represented by a
difference equation like that shown in Sect. 6.1. Consequently, the linear multi-step
methods are independent of whether the integrand function f(t, x) be, or be not, a
linear function of x.
Unlike the single-step methods, the multi-step methods require special tech-
niques, which will be described in the following paragraphs, to compute a sufficient
number of preceding values needed to start the integration process. The two prin-
ciples on which the quadrature formulae are based are, therefore: (a) the use of
preceding values of x and/or x′ to form a polynomial which approximates the
integrand function; and (b) the extrapolation of this polynomial into the next step of
the integration interval. The order k of the method is equal to the power of the step
size h in the total error term of the integration formula, which in turn is equal to one
more than the degree of the polynomial which interpolates the preceding values.
A multi-step (k-step) predictor-corrector method can in general be represented by
a difference equation of the following form
ak xi þ k þ ak1 xi þ k1 þ þ a0 xi ¼ hðbk fi þ k þ bk1 fi þ k1 þ þ b0 fi Þ
where fi = f(ti, xi), fi+1 = f(ti+1, xi+1), …, fi+k = f(ti+k, xi+k); {as} and {bs} (s = 0, 1,
…, k) are two sets of suitably chosen constants. All possible difference equations of
the form indicated above belong to two classes.
The first class, characterised by bk = 0, is such that the right-hand side of the
equation written above does not include fi+k, and therefore, the equation can be
solved explicitly for xi+k. The formulae of the first class (bk = 0) are called predictor
or open formulae. The second class, characterised by bk 6¼ 0, is such that the
right-hand side of the equation does include xi+k by means of the function fi+k, and
therefore, the equation cannot be solved explicitly for xi+k. In the latter case, the
solution is obtained by means of an iterative procedure and the corresponding
formulae are called corrector or closed formulae.
The constants {as} and {bs} (s = 0, 1, …, k) of a multi-step predictor-corrector
method are the counterparts of the parameters {ci}, {bi}, and {aij} (i, j = 1, 2, …, s)
characterising a Runge–Kutta formula (Sect. 6.2 et seq.) and, like them, cannot be
taken arbitrarily. On the contrary, they must be determined so that a given
multi-step formula should have the desired properties. One of these properties is the
accuracy of the method, which is measured by its local truncation error. The local
truncation error ei+1 of a multi-step method characterised by the following differ-
ence equation
ak xi þ k þ ak1 xi þ k1 þ þ a0 xi ¼ hðbk fi þ k þ bk1 fi þ k1 þ þ b0 fi Þ
is defined as follows
6.13 Multi-step Methods 951
X
k X
k
i þ 1 ¼ as xðti shÞ þ h bs f ðti shÞ xðti þ hÞ
s¼0 s¼1
which is the amount by which the numerical solution xi+1 differs from the exact
solution x(ti + h) x(ti+1) at the point t = ti+1.
As has been shown for the Runge–Kutta method, the expression of the local
truncation error can be expanded in a Taylor series in h. Following Dormand [24]
and remembering that x′(ti) = fi, this expansion may be written as follows
" #
X
k
ðshÞ2 00 ðshÞ3 ð3Þ
0
i þ 1 ¼ as xðti Þ shx ðti Þ þ x ðti Þ x ðti Þ þ
s¼0
2! 3!
" #
Xk
ðshÞ2 ð3Þ ðshÞ3 ð4Þ
0 00
þh bs x ðti Þ shx ðti Þ þ x ðt i Þ x ðt i Þ þ
s¼0
2! 3!

0 h2 00 h3 ð3Þ
xðti Þ þ hx ðti Þ þ x ðti Þ þ x ðti Þ þ
2! 3!
The terms can be so ordered as to obtain a series of powers of h, as follows

! !
X
k X
k X
k
0
i þ 1 ¼ xðti Þ þ as 1 þ hx ðti Þ bs sas 1 þ h2 x00 ðti Þ
s¼0 s¼1 s¼0
! !
X
k
1X k
1 1X k
1X k
1
3 ð3 Þ
sbs þ s as
2
þ h x ðti Þ s bs
2
s as
3
s¼1
2 s¼0 2 2 s¼1 3! s¼0 3!
!
1 X 3k
1X 4k
1
þ h4 xð4Þ ðti Þ s bs þ s as þ
3! s¼1 4! s¼0 4!
Assuming a−1 = 0, the term containing hp can be written as follows

( " # )
h p ð pÞ p1
X k
x ðti Þ ð1Þ s ðpbs sas Þ 1
p1
p! s¼1
In order for a multi-step predictor-corrector formula to be of the pth order, its

local truncation error must be O(hp+1). To this end, all terms of the local truncation
error up to the term containing hp must be equal to zero. Unlike the Runge–Kutta
formulae, here only one additional condition is necessary to increase the order of a
given formula from p to p + 1. The consistency conditions require that a
predictor-corrector formula should be at least of the first order.
Predictor-corrector methods are methods based on a predictor formula followed
by one or more applications of a corrector formula. As has been shown by Chase
[49], the predictor-corrector methods can be applied in more than one mode of
operation. In the first (or iterative) mode, the first estimate xpi+1 of the solution x at
the next mesh point i + 1 is followed by a corrector formula applied iteratively until
convergence is reached, as indicated below:
• an extrapolation formula is used to obtain a rough (or predicted) value of the
solution x at the end of the current step (this value is identified by the superscript
p, which stands for predictor: xpi+1);
• this predicted value is substituted into the differential equation, in order to
evaluate the integrand function as a function of the predicted value, i.e., fi+1 = f
(ti+1, xpi+1);
• fi+1 is used in a corrector formula to obtain a more accurate (or corrected) value
of xi+1 (this value is identified by the superscript c, which stands for corrector:
xci+1);
• this corrected value is substituted into the differential equation x′= f(t, x) and an
improved value of the integrand function is computed, as follows fi+1 = f(ti+1,
xci+1).
In the case of no further iterations performed by means of the corrector formula,
the process can now be advanced to the next step. By so doing, the number of
evaluations of the integrand function is equal to one plus the number of the cor-
rector iterations, and the stability of the method depends on the corrector formula
alone, because the predictor formula can only determine the number of iterations
required to reach convergence. In the second mode, the first estimate xpi+1 of the
solution is followed by only one application of the corrector formula, because the
predicted (xpi+1) and the corrected (xci+1) values are compared to obtain an estimate of
the truncation error committed in the current step. If this estimate does not exceed a
tolerance value specified by the solver, then the corrected value is accepted and is
taken as the final value for the current step; otherwise, the corrected value is
rejected and the step size is decreased before the current step is repeated starting
from the last accepted mesh point. By so doing, two evaluations of the integrand
function suffice to compute the solution at the next mesh point. Likewise, in case of
the error estimate being much smaller than the fixed tolerance, the step size may be
increased. A step-size change requires an interpolation of previous values or a
method restarting with a new (increased or decreased) step size. The second mode
of application of a predictor-corrector formula is much more used than the first; that
is, the need for successive recorrections is usually avoided by choosing a suffi-
ciently small step size. This second mode of application, in turn, may be used either
as described above or in a modified manner, as suggested by Hamming [50], who
applies the predictor equation followed by one application of the corrector equation,
and incorporates the error estimate xpi+1 − xci+1 in the final value xi+1 of the solution.
6.14 The Adams Method 953
6.14 The Adams Method
By integrating formally the differential equation x′= f(t, x), we have
Zti þ 1
xi þ 1 xi ¼ f ½t; xðtÞdt
ti
The integrand is the function f(t, x), which is replaced here by a polynomial p(t)
passing through some (e.g., four) points equispaced at the respective abscissae ti−3,
ti−2, ti−1, and ti. Let fi−3, fi−2, fi−1, and fi be the corresponding (known) ordinates of
the integrand function f(t, x). We want to extrapolate this function to the point
placed at the abscissa ti+1 = ti + h (predictor formula). By using a polynomial
p(t) of degree 3 or less, passing through the four points, to approximate the inte-
grand function f(t, x), we have (see [46]):
Zti þ 1
xi þ 1 xi pðtÞdt ¼ b0 f i3 þ b1 f i2 þ b2 f i1 þ b3 f i
ti
where the four constants b0, b1, b2, and b3 are to be determined.
When the origin of the t-axis is chosen to coincide with ti, i.e., ti = 0, then the
integral written above is taken over the interval going from 0 to h, where h is the
equal distance of the abscissae (h = ti+1 − ti).
Therefore, we have
Zh
f ðtÞdt b0 f ð3hÞ þ b1 f ð2hÞ þ b2 f ðhÞb3 f ð0Þ
0
The computation is performed below by replacing f(t) by successively t3, t2, t,

and 1. This yields
h4
¼ b0 ð3hÞ3 þ b1 ð2hÞ3 þ b2 ðhÞ3 þ b3 ð0Þ
4
h3
¼ b0 ð3hÞ2 þ b1 ð2hÞ2 þ b2 ðhÞ2 þ b3 ð0Þ
3
h2
¼ b0 ð3hÞ þ b1 ð2hÞ þ b2 ðhÞ þ b3 ð0Þ
2
h ¼ b0 ð1Þ þ b1 ð1Þ þ b2 ð1Þ þ b3 ð1Þ
Solving this system of algebraic equations for the unknowns b0, b1, b2, and b3,
we have
9 37 59 55
b0 ¼ h b1 ¼ h b2 ¼ h b3 ¼ h
24 24 24 24
The resulting fourth-order formula of the Adams-Bashforth predictor is then
1
xpiþ 1 ¼ xi þ hð55f i 59f i1 þ 37f i2 9f i3 Þ þ O h5
24
Now, having the predicted value xpi+1, we compute fi+1 = f(ti+1, xpi+1) and get the
improved (corrected) value xci+1 by interpolating the function f(t, x) between the
four points having co-ordinates (ti−2, fi−2), (ti−1, fi−1), (ti, fi) and (ti+1, fi+1).
The system of algebraic equations is
h4
¼ b 0 ð2hÞ3 þ b 1 ðhÞ3 þ b 2 ð0Þ3 þ b 3 ðhÞ3
4
h3
¼ b 0 ð2hÞ2 þ b 1 ðhÞ2 þ b 2 ð0Þ2 þ b 3 ðhÞ2
3
h2
¼ b 0 ð2hÞ þ b 1 ðhÞ þ b 2 ð0Þ þ b 3 ðhÞ
2
h ¼ b 0 ð1Þ þ b 1 ð1Þ þ b 2 ð1Þ þ b 3 ð1Þ
Solving this system for the unknowns b*0, b*1, b*2, and b*3, we have
1 5 19 9
b 0 ¼ h b 1 ¼ h b 2 ¼ h b 3 ¼ h
24 24 24 24
The resulting fourth-order formula of the Adams-Moulton corrector is then
1
xciþ 1 ¼ xi þ h 9f i þ 1 þ 19f i 5f i1 þ f i2 þ O h5
24
If the predicted and corrected values agree within the prescribed tolerance, the
step size may be increased.
A possible way for doing this is to double the step size. In the fourth-order
formulae considered above, we can do this after computing seven values of x at
equidistant abscissae by dropping every second one. Of course, it is also possible to
restart the multi-step method, using any desired step size, by means of a single-step
method. Inversely, if the difference between predicted and corrected values is greater
than the prescribed tolerance, the step size must be decreased, as the sequel will show.
The Adams predictor-corrector formulae written above have been derived using
four points and a polynomial of degree three or less. However, the same principle
may be applied to derive Adams formulae of higher order.
For example, in case of five points, the coefficients for the predictor formula are
251 1274 2616 2774 1901

b0 ¼ h b1 ¼ h b2 ¼ h b3 ¼ h b4 ¼ h
720 720 720 720 720
Therefore, the fifth-order formula of the Adams-Bashforth predictor is
1
xpiþ 1 ¼ xi þ hð1901f i 2774f i1 þ 2616f i2 1274f i3 þ 251f i4 Þ þ O h6
720
Likewise, in case of five points, the coefficients for the corrector formula are
19 106 264 646 251

b 0 ¼ h b 1 ¼ h b 2 ¼ h b 3 ¼ h b 4 ¼ h
720 720 720 720 720
Therefore, the fifth-order formula of the Adams-Moulton corrector is
1
xciþ 1 ¼ xi þ h 251f i þ 1 þ 646f i 264f i1 þ 106f i2 19f i3 þ O h6
720
In the general case of a i-step Adams method, we use a polynomial of degree less
than or equal to i expressed in terms of backward Newton-Gregory differences.
For this purpose, we introduce the variable
t ti
s¼
h
The derivative x′= f(t, x) of the unknown function x(t) is replaced by a poly-
nomial p(t) having degree less than or equal to i and x′ is integrated over the desired
interval. The appropriate polynomial p(t) is given by Newton’s formula for back-
ward interpolation
s ð s þ 1Þ 2 s ð s þ 1Þ ð s þ 2Þ 3
pi ðtÞ ¼pi ðti þ hsÞ ¼ xi þ srxi þ r xi þ r xi þ
2! 3!
s ð s þ 1Þ ð s þ 2Þ ð s þ i 1Þ i
þ r xi
i!
where
rxk ¼ xk xk1 r2 xk ¼ rðrxk Þ . . . ri xk ¼ rðri1 xk Þ
are the backward differences of xk.

Writing the interpolation formula for the derivative (x′) leads to
sðs þ 1Þ 2 0
pi ðtÞ ¼ pi ðti þ hsÞ ¼ x0 i þ srx0 i þ r xi
2!
s ð s þ 1Þ ð s þ 2Þ 3 0 sðs þ 1Þðs þ 2Þ ðs þ i 1Þ i 0
þ r x i þ þ rxi
3! i!
1 1
¼ x0 i þ srx0 i þ s2 þ s r2 x0 i þ s3 þ 3s2 þ 2s r3 x0 i
2 6
1 4
þ s þ 6s3 þ 11s2 þ 6s r4 x0 i þ
24
The change of x for any desired interval results from integrating x′ over the same
interval, as follows
Ztk þ 1
Dx ¼ x0 dt
tk
By substituting pi(t) for x′ in the preceding expression, we have

Ztk þ 1
1 2 1
Dx ¼ x0 i þ srx0 i þ
s þ s r2 x0 i þ s3 þ 3s2 þ 2s r3 x0 i
2 6
tk

1 4
þ s þ 6s3 þ 11s2 þ 6s r4 x0 i þ . . . dt
24
From the definition s = (t − ti)/h, it follows that

hs ¼ t ti
t ¼ ti þ hs
dt ¼ h ds
Substituting the value hds for dt in the preceding expression of Dx and changing
limits of integration, we have
Zsk þ 1
1 2 1
Dx ¼ h x0 i þ srx0 i þ
s þ s r2 x0 i þ s3 þ 3s2 þ 2s r3 x0 i
2 6
sk

1 4 4 0
þ s þ 6s þ 11s þ 6s r x i þ ds
3 2
24
2

0 0 s 1 2 0 s3 s2 1 3 0 s4
¼ x i s þ rx i þ r x i þ þ r xi þs þs
3 2
2 2 3 2 6 4
5 sk þ 1
1 4 0 s s4 s3
þ r xi þ 3 þ 11 þ 3s2 þ
24 5 2 3 sk
Now, let us compute the value of Dx for the intervals [ti, ti+1] and [ti−1, ti] by
substituting the proper limits for s in the preceding expression.
For the interval [ti, ti+1], the limits of s are
ti þ 1 ti h ti ti
sk þ 1 ¼ ¼ ¼1 sk ¼ ¼0
h h h
Hence

1 5 2 0 3 251 4 0
xpiþ 1 xi ¼ ½Dxii þ 1 ¼ h x0 i þ rx0 i þ r x i þ r 3 x0 i þ r xi þ
2 12 8 720
where the superscript p indicates the predictor formula in difference form.

For the interval [ti−1, ti], the limits of s are
ti ti
sk þ 1 ¼ ¼0
h
ti1 ti h
sk ¼ ¼ ¼ 1
h h
Hence

0 1 0 1 2 0 1 3 0 19 4 0
xi xi1 ¼ ½Dxii1 ¼ h x i rx i r x i r x i r x iþ
2 12 24 720
If we advance one step, the same formula is

1 1 1 19 4 0
xciþ 1 ¼ xi þ h x0 i þ 1 rx0 i þ 1 r2 x0 i þ 1 r3 x0 i þ 1 r x iþ1
2 12 24 720
where the superscript c indicates the corrector formula in difference form.

As an example of application (from [51]), let us compute the numerical solution
of the following differential equation x′= t + x, with the initial condition t0 = 0,
x0 = 1, by means of the Adams fourth-order formulae, by using a step size
h = 0.05. Since the method is not self-starting (see Sect. 6.13), at least four points
must be known in addition to the initial point (t0, x0). Let us suppose that all of the
following points are known
t2 ¼ 0:10 x2 ¼ 0:9096748

t1 ¼ 0:05 x1 ¼ 0:9524588
t0 ¼ 0:00 x0 ¼ 1:000000
t1 ¼ 0:05 x1 ¼ 1:052542
t2 ¼ 0:10 x2 ¼ 1:110342
These values may be computed by means of a single-step method, for example,

by means of the Runge–Kutta method applied forward and backward with respect
to the initial point (t0, x0), as will be shown in Sect. 6.17. By means of these values,
we can construct the following table.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

−0.10 0.9096748 0.8096748
−0.05 0.9524588 0.9024588 0.0927840
0.00 1.0000000 1.0000000 0.0975412 0.0047572
0.05 1.0525420 1.1025420 0.1025420 0.0050008 0.0002436
0.10 1.1103420 1.2103420 0.1078000 0.0052580 0.0002572 0.0000136
0.15
The first column (t) is filled by the abscissae t−2, t−1, t0, t1, and t2 equispaced by
the step size h = 0.05. The second column (x) is filled by the corresponding
ordinates x−2, x−1, x0, x1, and x2. The third column (x′) is filled by evaluating the
integrand function x′−2 = t−2 + x−2, x′−1 = t−1 + x−1, and so on. Each of the fol-
lowing columns is filled by computing, respectively, the first, the second, the third,
and the fourth backward difference of x′, as follows
rx0 0:05 ¼ x0 0:05 x0 0:10 ¼ 0:9024588 0:8096748 ¼ 0:0927840

rx0 0:00 ¼ x0 0:00 x0 0:05 ¼ 1:0000000 0:9024588 ¼ 0:0975412
..
.
r2 x0 0:00 ¼ rx0 0:00 rx0 0:05 ¼ 0:0975412 0:0927840 ¼ 0:0047572
..
.
In order to compute the values for the next row (i + 1), corresponding to t = 0.15,
we apply the fourth-order predictor formula in difference form, as follows

1 5 2 0 3 251 4 0
xpiþ 1 ¼ xi þ h x0 i þ rx0 i þ r x i þ r3 x0 i þ r xi
2 12 8 720

1 5
¼ 1:110342 þ 0:05 1:210342 þ 1:1078 þ 0:005258
2 12

3 251
þ 0:0002572 þ 0:0000136 ¼ 1:173669
8 720
where the index of row (i) points to the row x = 0.10. Hence, xi = x0.10 = 1.110342
and xi+1 = x0. Using the value given by the predictor, we can compute x′i+1, as
follows

x0 i þ 1 ¼ f ti þ 1 ; xpiþ 1 ¼ ti þ 1 þ xpiþ 1 ¼ 0:15 þ 1:173669 ¼ 1:323669
Therefore, we can fill the next row (t = 0.15) of the table as follows.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

−0.10 0.9096748 0.8096748
−0.05 0.9524588 0.9024588 0.0927840
0.00 1.0000000 1.0000000 0.0975412 0.0047572
0.05 1.0525420 1.1025420 0.1025420 0.0050008 0.0002436
0.10 1.1103420 1.2103420 0.1078000 0.0052580 0.0002572 0.0000136
0.15 1.3236690 0.1133270 0.0055270 0.0002690 0.0000118
Now, we apply the fourth-order corrector formula in difference form, as follows

1 1 1 19 4 0
xciþ 1 ¼ xi þ h x0 i þ 1 rx0 i þ 1 r2 x0 i þ 1 r3 x0 i þ 1 r x iþ1
2 12 24 720

1 1
¼ 1:110342 þ 0:05 1:323669 0:113327 0:005527
2 12

1 19
0:000269 0:0000118 ¼ 1:173669
24 720
This value is the same as that which has been computed above by means of the
predictor formula, and therefore, it is accepted as the value of x0.15.
Now, the index of row (i) is increased by one, that is, i points to the row t = 0.
We are ready to perform the next step in the same way as that shown above.
In the general case, if we indicate with cm the constants computed above for
predicted value of xi+1, we can express them (see [2]) as follows
Zti þ 1 Z1
m1 s s
cm ¼ ð1Þ dt ¼ ð1Þm ds
h m m
ti 0
The values of such constants do not depend on the integrand function f(t, x) and
are given in the following table.
m 0 1 2 3 4 5 6
cm 1 1/2 5/12 3/8 251/720 95/288 19087/60480
The predictor constants can also be computed by means of the following

formula:
1 1 1
cm þ cm1 þ cm2 þ þ c ¼1
2 3 mþ1 0
Therefore, we have
c0 ¼ 1
1 1 1
c1 ¼ 1 c0 ¼ 1 1 ¼
2 2 2
1 1 1 1 1 5
c2 ¼ 1 c1 c0 ¼ 1 1 ¼
2 3 2 2 3 12
..
.
When the differences are expressed in terms of ordinates, the Adams-Bashforth

formula is
X
k
xpiþ 1 xi ¼ h bkn f in
n¼0
where the coefficients bkn are given by

n nþ1 k
bkn ¼ ð1Þn cn þ cn þ 1 þ þ c ðn ¼ 0; 1; . . .; k; k ¼ 0; 1; . . .Þ
n n n k
The values of bkn are given in the following table.
n 0 1 2 3 4 5
b0n 1
b1n 3/2 −1/2
b2n 23/12 −16/12 5/12
b3n 55/24 −59/24 37/24 −9/24
b4n 1901/720 −2774/720 2616/720 −1274/720 251/720
b5n 4227/1440 −7673/1440 9482/1440 −6798/1440 2627/1440 −425/1440
In the same manner, the corrector constants c*m computed above
Zti Z0
m1 s m s
c m ¼ ð1Þ dt ¼ ð1Þ ds
h m m
ti1 1
do not depend on f(t, x) and are given in the following table.
m 0 1 2 3 4 5 6
c*m 1 −1/2 −1/12 −1/24 −19/720 −3/160 −863/60480
Such constants can also be computed by means of the following formulae:
1 1 1
c m þ c m1 þ c m2 þ þ c ¼ 1 ð m ¼ 0Þ
2 3 m þ 1 m0
1 1 1
c m þ c þ c þ þ c ¼ 0 ðm ¼ 1; 2; . . .Þ
2 m1 3 m2 m þ 1 m0
Therefore, we have
c 0 ¼ 1
1 1 1
c 1 ¼ c 0 ¼ 1 ¼
2 2 2
1 1 1 1 1 1 1
c 2 ¼ c 1 c 0 ¼ 1¼ 1¼
2 3 2 2 2 3 12
..
.
When the differences are expressed in terms of ordinates, the Adams-Moulton

formula is
X
k
xciþ 1 xi ¼ h b kn f in
n¼0
where the coefficients b*kn are given by

n nþ1 k
b kn ¼ ð1Þn c nþ c nþ1 þ þ ck
n n n
ðn ¼ 0; 1; . . .; k; k ¼ 0; 1; . . .Þ
The values of b*kn are given in the following table.
n 0 1 2 3 4 5
b*0n 1
b*1n 1/2 1/2
b*2n 5/12 8/12 −1/12
b*3n 9/24 19/24 −5/24 1/24
b*4n 251/720 646/720 −264/720 106/720 −19/720
b*5n 475/1440 1427/1440 −798/1440 482/1440 −173/1440 27/1440
The formulae given above may also be derived (see [52]) by using the so-called
linear operators, some of which have been shown above. They are
Df i ¼ f i þ 1 f i ðfirstÞ forward difference

rf i ¼ f i f i1 ðfirstÞ backward difference
Df i ¼ f i þ 1=2 f i1=2 ðfirstÞ central difference
As has been shown above, further differences than the first (i.e., the second, the
third, …, the nth) are defined as follows
r2 f i ¼ rrf i ¼ rðf i f i1 Þ ¼ f i f i1 ðf i1 f i2 Þ ¼ f i 2f i1 þ f i2

..
.
rn f i ¼ rn1 f i rn1 f i1
The displacement operator E is defined as follows

Ef i ¼ f i þ 1
and its second, third, …, nth powers are defined as follows
E2 f i ¼ f i þ 2 E3 f i ¼ f i þ 3 ... En f i ¼ f i þ n
The definition given above (∇fi = fi − fi−1) of the first backward difference operator
shows that this operator may be expressed in terms of the displacement operator E
r ¼ 1 E1
Conversely, the displacement operator may be expressed in terms of the back-

ward difference operator ∇
1
E¼
1r
Negative powers of the backward difference operator ∇ (i.e., ∇−1, ∇−2) indicate
summation, not differentiation. Therefore, remembering the definition
rf i ¼ f i f i1
and applying ∇−1 to both sides, we have
r1 ðrf i Þ ¼ r1 ðf i f i1 Þ
that is,
f i ¼ r1 f i r1 f i1

Hence, the summation operator (∇−1) is defined by means of the recurrence
r1 f i ¼ f i r1 f i1
The initial term (corresponding to i = 0) in the sum is arbitrary, this term being an
integration constant (c1). Therefore, we have
r1 f i ¼ c1 þ f 1 þ f 2 þ þ f i ði
1Þ
r1 f i ¼ c1 f i þ 1 f i þ 2 f 0 ði 1Þ
The second power (∇−2) of the summation operator corresponds to the second
power (∇2) of the difference operators and appears in the Gauss-Jackson formula (as
the sequel will show), also called second-sum integration formula. The second sum
(∇−2) has an integration constant c2, analogous to c1.
The differentiation operator D is defined as the operator that, when applied to a
function f, turns f into the derivative of f with respect to the independent variable:
dð. . .Þ
Dð. . .Þ ¼
dt
The differentiation operator may be expressed in terms of the displacement oper-

ator, as follows
ðphÞ2 2
E p f ðt0 Þ ¼ f ðt0 þ phÞ ¼ f ðt0 Þ þ ph D f ðt0 Þ þ D f ðt0 Þ
2!
ðphÞ3 3
þ D f ðt0 Þ þ ¼ Eph D f ðt0 Þ
3!
where p is any real number. The exponential of an operator indicates the Taylor
expansion of the operator, as follows
1 2 1 3
et expðtÞ ¼ 1 þ t þ t þ t þ
2! 3!
Since
E p f ðt0 Þ ¼ EphD f ðt0 Þ
it follows that
E p ¼ ephD
or
phD ¼ log E p ¼ p log E
The indefinite integration operator D−1 is defined as the operator inverse of dif-
ferentiation, as follows
Z
ð. . .Þdt ¼ D1 ð. . .Þ
In other words, the indefinite integration operator, applied to an arbitrary function f,

produces the primitive (i.e., the anti-derivative) of that function. As has been shown
above

log E 1 1 logð1 rÞ
D¼ ¼ log ¼
h h 1r h
the single-integration operator D−1 may, therefore, be expressed as the inverse of

the differentiation operator D, as follows
h
D1 ¼
logð1 rÞ
The definite integral operator is computed by applying the displacement operator

E to the indefinite integral. The following discussion takes in consideration, without
loss of generality, the second-order differential equation
x0 ¼ y
y0 ¼ f ðt; xÞ
where x and y are, respectively, the position vector and the velocity vector
describing the motion of a particle, f(t, x) is the integrand function, and the prime
sign denotes first derivatives with respect to time (i.e., x′ dx/dt and y′ dy/dt),
which is the independent variable. The position and velocity vectors must be
specified at the time t = t0, that is, x(t0) = x0 and y(t0) = y0.
The definite integral of f(t, x) may be expressed as follows
tZ
i þ ph
1
ðE 1ÞD f i ¼ ðE 1Þyi ¼ ðE 1Þ
p p p
f ðt; xÞdt
ti
where p is equal to 1 if the integration is performed over one time step. The
expression written above is a predictor formula, because a value of the velocity
vector y is computed by advancing one step (h) from the current abscissa (ti) to the
next abscissa (ti + h). The predictor operator J is defined as follows
h i rD1
J ¼ h1 ðE 1ÞD1 ¼ h1 ð1 rÞ1 1 D1 ¼
hð 1 r Þ
r
¼
ð1 rÞ logð1 rÞ
The corrector operator L results from integrating from ti − h to ti. To this end, the
expression given above is shifted backward by using the displacement operator:
r
L ¼ E1 J ¼ ð1 rÞJ ¼
logð1 rÞ
The coefficients of the operators J and L may be computed by using a recursion

relation. Let us consider, in the first place, the corrector operator L, which may be
expressed by its Taylor-series expansion using coefficients c*i, as follows
r
L ¼ E 1 J ¼ ¼ c 0 þ c 1 r þ c 2 r 2 þ c 3 r 3 þ
logð1 rÞ
Since the function log(1 − ∇), expanded in Taylor series, is
1 1 1
logð1 rÞ ¼ r r2 r3 r4
2 3 4
then we have
r 1 X1
L¼ ¼ ¼ c i r i
logð1 rÞ 1 þ 2 r þ 3 r2 þ 4 r3 þ
1 1 1
i¼0
Hence,

1
1 2 1 3
c 0 þ c 1r þ c 2r þ c 3r þ 1 þ r þ r þ r þ ¼ 1
2 3
2 3 4
Expanding the product and grouping by powers of ∇, we have

1 1 1 1 1 1
c 0 þ c 0 þ c 1 r þ c 0 þ c 1 þ c 2 r2 þ c 0 þ c 1 þ c 2 þ c 3 r3 þ ¼ 1
2 3 2 4 3 2
It follows that
c 0 ¼ 1
1
c þ c 1 ¼ 0 hence c 1 ¼ 12
2 0
1 1
c 0 þ c 1 þ c 2 ¼ 0 hence c 2 ¼ 12
1
3 2
1 1 1
c 0 þ c 1 þ c 2 þ c 3 ¼ 0 hence c 3 ¼ 24
1
4 3 2
..
.
X
n1
c i
c n ¼ ðn
1; c 0 ¼ 1Þ
i¼0
nþ1 i
It is evident that the coefficients c*0, c*1, c*2, c*3, … are the same as those which
have been computed above for the Adams-Moulton corrector formula.
Once the values of the coefficients c*i are known, the integration over the time
interval [ti − h, ti] is performed as follows
Zti
f ðt; xÞdt ¼ hE1 J f i ¼ h L f i
ti h
and the corrected value is

1 1 1
yi ¼ yi1 þ h 1 r r2 r3 þ f i
2 12 24
that is, the Adams-Moulton formula in difference form which has been derived
above.
In order to derive the predictor formula by means of the linear operators, we note
(see above) that
J ¼ ð1 rÞ1 L
The Taylor-series expansion of (1 − ∇)−1 is 1 + ∇ + ∇2 + ∇3 + and the

Taylor-series expansion of L has been found above in terms of the coefficients c*i.
Therefore, we have

J ¼ ð1 rÞ1 L ¼ 1 þ r þ r2 þ r3 þ c 0 þ c 1 r þ c 2 r2 þ c 3 r3 þ
This equation may be written as a Taylor-series expansion in terms of coefficients

c0, c1, c2, …, such that
ð1 þ r þ r2 þ r3 þ Þðc 0 þ c 1 r þ c 2 r2 þ c 3 r3 þ Þ
¼ c0 þ c1 r þ c2 r 2 þ c3 r 3 þ
Expanding the product and grouping by powers of ∇, we have
c 0 þ ðc 0 þ c 1 Þr þ ðc 0 þ c 1 þ c 2 Þr2 þ ¼ c0 þ c1 r þ c2 r2 þ c3 r3 þ
Hence
c0 ¼ c 0 ¼ 1
1 1
c1 ¼ c 0 þ c 1 ¼ 1 ¼
2 2
1 1 5
c2 ¼ c 0 þ c 1 þ c 2 ¼ 1 ¼
2 12 12
1 1 1 3
c3 ¼ c 0 þ c 1 þ c 2 þ c 3 ¼ 1 ¼
2 12 24 8
..
.
Therefore, each coefficient ci is the sum of the coefficients c*k, where k i, that is,
X
i
ci ¼ c k
k¼0
The coefficients c0, c1, c2, c3, … are evidently the same as those which have been
computed above for the Adams-Bashforth predictor formula.
Now, the integration of f(t, x) over the time interval [ti, ti + h] is performed as
follows
Z
ti þ h
f ðt; xÞdt ¼ h J f i
ti
and the predicted value is

1 5 2 3 3
yi þ 1 ¼ yi þ h 1 þ r þ r þ r þ fi
2 12 8
that is, the Adams-Bashforth formula in difference form which has been derived
above. Summarising, we have derived the Adams predictor-corrector formulae
firstly in ordinate form (coefficients bi and b*i found directly), secondly in differ-
ence form (coefficients ci and c*i found by means of interpolating polynomials
expressed in terms of backward Newton-Gregory differences), and thirdly in dif-
ference form by means of linear operators. As shown above, it is immediate to
convert the difference form into the ordinate form. In the following section of this
paragraph, another version, called the summed form, of the Adams formulae will be
derived by means of the summation operator (∇−1).
To this end, we start from the Adams-Moulton corrector formula (see above)

1 1 1
yi yi1 ¼ ryi ¼ h 1 r r2 r3 þ f i
2 12 24
By multiplying both sides of the preceding equation by ∇−1, we have

1 1 1
yi ¼ h r1 r r2 þ f i
2 12 24
which is the summed form of the Adams-Moulton corrector formula.

In the same way, we start from the Adams-Bashforth predictor formula (see
above)

1 5 2 3 3
yi þ 1 yi ¼ ryi þ 1 ¼ h 1þ rþ r þ r þ fi
2 12 8
which, multiplied by ∇−1, yields

1 1 5 3 2
yi þ 1 ¼h r þ þ r þ r þ fi
2 12 8
which is the summed form of the Adams-Bashforth predictor formula.
6.15 The Störmer-Cowell Method
This method applies to the special class of second-order differential equations

considered in Sects. 6.7 and 6.8, to discuss the Runge–Kutta–Nyström method of
integration. Let us consider again the special second-order equation x″= f(t, x) written
in vector form as follows
x0 ¼ y
y0 ¼ f ðt; xÞ
where x and y are, respectively, the position vector and the velocity vector
describing the motion of a particle, f(t, x) is the integrand function, and the prime
sign denotes first derivatives with respect to time (i.e., x′ dx/dt and y′ dy/dt),
which is the independent variable. The position and velocity vectors must be
specified at the time t = t0, that is, x(t0) = x0 and y(t0) = y0. The integrand function
f(t, x) does not depend explicitly on y. The step size h denotes the time interval
starting from the initial time t = t0.
6.15 The Störmer-Cowell Method 969
The Adams-Moulton corrector formula derived in Sect. 6.14
Zti
f ðt; xÞdt ¼ h E1 J f i ¼ h L f i
ti h
is the corrector formula for the velocity vector y.

In order to find the corresponding formula for the position vector x, we multiply
the preceding expression by the corrector operator, as follows
Zti
hL f ðt; xÞdt ¼ h Lðyi yi1 Þ ¼ h2 L2 f i
ti h
The corrector operator, applied to the velocity vector yi, yields the expression
hLyi ¼ xi xi1
It follows that
hLðyi yi1 Þ ¼ ðxi xi1 Þ ðxi1 xi2 Þ ¼ xi 2xi1 þ xi2 ¼ h2 L2 f i
That is,
xi ¼ 2xi1 xi2 þ h2 L2 f i
The operator L2 may be expanded in a Taylor series as follows
L2 ¼ r 0 þ r 1 r þ r 2 r2 þ r 3 r3 þ
Now, remembering the expression of Sect. 6.14
L ¼ c 0 þ c 1 r þ c 2 r 2 þ c 3 r 3 þ
we can write
L2 ¼ r 0 þ r 1 r þ r 2 r2 þ r 3 r3 þ
2
¼ c 0 þ c 1 r þ c 2 r 2 þ c 3 r 3 þ
Expanding and grouping the powers of ∇, we have:
r 0 ¼ c 0 c 0
r 1 ¼ c 0 c 1 þ c 1 c 0
r 2 ¼ c 0 c 2 þ c 1 c 1 þ c 2 c 0
r 3 ¼ c 0 c 3 þ c 1 c 2 þ c 2 c 1 þ c 3 c 0
..
.
X
i
r i ¼ c k c ik
k¼0
The values of c*k have been computed in Sect. 6.14. They are
k 0 1 2 3 4 5 6
c*k 1 −1/2 −1/12 −1/24 −19/720 −3/160 −863/60480
Therefore, the values of the coefficients r*i, computed by means of the coeffi-
cients c*k, are
r 0 ¼ 1 1 ¼ 1

1 1
r 1 ¼1 þ 1 ¼ 1
2 2

1 1 1 1 1
r 1 ¼1 þ þ 1¼
12 2 2 12 12
..
.
that is,
k 0 1 2 3 4 5 6
r*k 1 −1 1/12 0 −1/240 −1/240 −221/60480
Using coefficients r*k in the equation written above
we have the Cowell corrector formula

xi ¼ 2xi1 xi2 þ h2 r 0 þ r 1 r þ r 2 r2 þ r 3 r3 þ f i

1 2 1 4 1 5
¼ 2xi1 xi2 þ h2 1 r þ r r r fi
12 240 240
As has been done in Sect. 6.14, we derive the predictor formula by simply mul-
tiplying both sides of the equation
by the shifting operator (E), that is,

Exi ¼ E 2xi1 xi2 þ h2 L2 f i
which results in
xi þ 1 ¼ 2xi xi1 þ h2 EL2 f i
Let us observe that (see Sect. 6.14)

EL2 ¼ ð1 rÞ1 L2 ¼ 1 þ r þ r2 þ r3 þ r 0 þ r 1 r þ r 2 r2 þ r 3 r3 þ
X1
¼ ri ri
i¼1
Expanding, grouping by powers of ∇, and using the values of r*k shown above,
we have
r0 ¼ r 0 ¼ 1
r1 ¼ r 1 þ r 0 ¼ 1 þ 1 ¼ 0
1 1
r2 ¼ r 2 þ r 1 þ r 0 ¼ 1þ1 ¼
12 12
1 1
r3 ¼ r 3 þ r 2 þ r 1 þ r 0 ¼ 0 þ 1þ1 ¼
12 12
..
.
X
i
ri ¼ r i þ r i1 þ þ r 1 þ r 0 ¼ r k
k¼0
The values of the coefficients rk, computed by means of the r*k, are given in the
following table.
k 0 1 2 3 4 5 6
rk 1 0 1/12 1/12 19/240 3/40 863/12096
By using the coefficients rk in the predictor expression written above
xi þ 1 ¼ 2xi xi1 þ h2 EL2 f i

we have the Störmer predictor formula

xi þ 1 ¼ 2xi xi1 þ h2 r0 þ r1 r þ r2 r2 þ r3 r3 þ f i

1 2 1 3 19 4 3 5
¼ 2xi xi1 þ h 1 þ
2
r þ r þ r þ r fi
12 12 240 40
An example of application of these formulae is given below. Let us consider the

following scalar differential equation of the second order
x00 ¼ x
defined for all values of the independent variable t with the initial conditions x0 = 1
and x′0 = −1 for t = 0. As is well known, the analytical solution of this initial-value
problem is x(t) = exp(−t). We want to integrate numerically this differential equa-
tion, with the specified initial conditions, by means of a Störmer-Cowell fourth-order
formula, with a step size h = 0.05. Let the following starting values be known:
t2 ¼ 0:10 x2 ¼ 1:1051709

t1 ¼ 0:05 x1 ¼ 1:0512711
t0 ¼ 0:00 x0 ¼ 1:0000000
t1 ¼ 0:05 x1 ¼ 0:9512294
t2 ¼ 0:10 x2 ¼ 0:9048374
For this purpose, we use the integration table shown below.
i ti xi fi ∇fi ∇2fi ∇3fi ∇4fi

−2 −0.10 1.1051709
−1 −0.05 1.0512711
0 0.00 1.0000000
1 0.05 0.9512294
2 0.10 0.9048374
To fill the table, we observe that the integrand function f(x), in the differential
equation given above, is such that fi = xi and thus we can fill the column of fi in the
table, by writing there the same values as those contained in the column of xi. The
result is given in the following table.

−2 −0.10 1.1051709 1.1051709
−1 −0.05 1.0512711 1.0512711
0 0.00 1.0000000 1.0000000
1 0.50 0.9512294 0.9512294
2 0.10 0.9048374 0.9048374
Now, operating as has been shown in Sect. 6.14, we can fill the columns of the
first, second, third, and fourth differences by inserting in each column the difference
of the values placed in the column on the left-hand side of the column considered,
as follows.

−2 −0.10 1.1051709 1.1051709
−1 −0.05 1.0512711 1.0512711 −0.0538998
0 0.00 1.0000000 1.0000000 −0.0512711 0.0026287
1 0.05 0.9512294 0.9512294 −0.0487706 0.0025005 −0.0001282
2 0.10 0.9048374 0.9048374 −0.0463920 0.0023786 −0.0001219 0.0000063
We are now ready to perform the first step of numerical integration by applying
the formulae given above.
Störmer predictor:

ð pÞ 1 2 1 3 19 4
x3 ¼ 2x2 x1 þ h2 f2 þ r f2 þ r f2 þ r f2
12 12 240

0:0023786
¼ 2 0:9048374 0:9512294 þ 0:052 0:9048374 þ
12

0:0001219 19
þ 0:0000063 ¼ 0:8607080
12 240
Evaluation of the integrand function:

ðpÞ
f3 ¼ f t3 ; x3 ¼ 0:8607080
Having computed f3 = 0.8607080, we can now use the values contained in the
row i = 2 of the integration table shown above to compute the differences
rf3 ¼ f3 f2 ¼ 0:8607080 0:9048374 ¼ 0:0441294

r2 f3 ¼ rf3 rf2 ¼ 0:0441294 ð0:0463920Þ ¼ 0:0022626
r3 f3 ¼ r2 f3 r2 f2 ¼ 0:0022626 0:0023786 ¼ 0:0001160
r4 f3 ¼ r3 f3 r3 f2 ¼ 0:0001160 ð0:0001219Þ ¼ 0:0000059
We insert f3 and these differences in the integration table, as follows.

−2 −0.10 1.1051709 1.1051709
−1 −0.05 1.0512711 1.0512711 −0.0538998
0 0.00 1.0000000 1.0000000 −0.0512711 0.0026287
(continued)
(continued)
1 0.05 0.9512294 0.9512294 −0.0487706 0.0025005 −0.0001282
2 0.10 0.9048374 0.9048374 −0.0463920 0.0023786 −0.0001219 0.0000063
3 0.15 0.8607080 −0.0441294 0.0022626 −0.0001160 0.0000059
Cowell corrector:

ðcÞ 1 2 1 4
x3 ¼ 2x2 x1 þ h2 f3 rf3 þ r f3 r f3
12 240

0:0022626
¼ 2 0:9048374 0:9512294 þ 0:052 0:8607080 þ 0:0441294 þ
12

0:0000059
¼ 0:8607080
240
This value is the same as that which has previously been found by means of the
predictor formula, and therefore, we can proceed to the next step. The exact
solution, with seven decimal figures, is x = exp(−0.15) = 0.8607080.
6.16 The Gauss-Jackson Method
The method shown in the present paragraph bears this name because it was orig-
inally found by Gauss and subsequently described by Jackson [53].
The formulae presented below are the summed form of the Störmer-Cowell
formulae derived in Sect. 6.15. They are also called second-sum formulae. Let us
consider again the Cowell corrector (see Sect. 6.15):

1 2 1 4 1 5
xi 2xi1 þ xi2 ¼ h2 1 r þ r r r fi
12 240 240
Observing that the left-hand side of this equation is the second backward difference
of the position vector xi, we have

1 2 1 4 1 5
r xi ¼ h
2 2
1 $þ $ $ $ fi
12 240 240
By applying the second-sum operator (∇−2) to both sides of this equation, we have

2 1 2 1 4 1 5
xi ¼ h r 1 rþ
2
r r r fi
12 240 240
6.16 The Gauss-Jackson Method 975
Since (1 − ∇) = E−1, it follows that (1 − ∇)fi = E−1fi and hence

1 1 2 1 3
xi ¼ h2 r2 f i1 þ r r fi
12 240 240
By using the Störmer predictor (see Sect. 6.15)

1 2 1 3 19 4 3 5
r 2 xi þ 1 ¼ h2 1 þ r þ r þ r þ r þ fi
12 12 240 40
and operating likewise, we have

12 1 19 2 3 3
xi þ 1 ¼ h r fi þ
2
þ rþ r þ r þ fi
12 12 240 40
which is the Gauss-Jackson predictor formula.

In the predictor formula, the second-sum operator (∇−2) is applied to fi, whereas,
in the corrector formula, the same operator is applied to fi−1.
In order to compute the negative powers of ∇ (i.e., ∇−1 and ∇−2), the backward
difference tables used in Sect. 6.15 must be extended to the left. This extension
yields
r1 ðrf i Þ ¼ f i ¼ r1 ðf i f i1 Þ ¼ r1 f i r1 f i1
This expression, in turn, yields a recursion formula which defines the inverse
operator ∇−1fi as follows
r1 f i ¼ f i þ r1 f i1
Thus, the values of the negative powers of ∇ can be introduced into the difference
table, as the following example will show.
The second power ∇−2 of the summation operator ∇−1 is a second sum, which
has its counterpart in the second difference ∇2.
As Dormand [24] points out, this method requires the initialisation of the first
and second sums, which are such that
r1 f i ¼ r2 f i r2 f i1

f i ¼ r1 f i r1 f i1
To this end, Merson [54] has derived the following formulae
1 1 1 11 3 11 4
r1 f i1 ¼ h1 x0 i þ f þ rf r2 f i þ 1 r fiþ2 þ r fiþ2 þ
2 i 12 i þ 1 24 720 1440
1 1 2 31
r2 f i ¼ h2 xi f i þ r fiþ1 r4 f i þ 2 þ
12 240 60480
The value of i is in the centre of the difference table, which is extended on its
left-hand side to include the values of the first and second sums.
Instead of using Merson’s formulae, it is possible of course to compute the
second sum by means of the Gauss-Jackson corrector formula given above, that is,

1 1 2 1 3 221 4
xi ¼ h2 r2 f i1 þ r r r fi
12 240 240 60480
This equation, solved for ∇−2fi−1, yields the second sum as follows

xi 1 1 2 1 3 221 4
r2 f i1 ¼ r r r fi
h2 12 240 240 60480
which is just the expression given by Dormand [24, p. 268] as follows
xi X k2
p1
r2 f i1 ¼ r r fi
h2 k¼2 k
where r*k are the Cowell corrector coefficients of Sect. 6.15, that is,
k 0 1 2 3 4 5 6
r*k 1 −1 1/12 0 −1/240 −1/240 −221/60480
and p is the order of the formula. An expression like this can be written for ∇−2fi−2.
By operating in this way, from a set of given starting values xn it possible to
compute the second-sum values necessary to fill the column of ∇−2fi in the inte-
gration table. The first-sum values result from the backward differences of the
second-sum values. This makes it possible to fill the column of ∇−1fi in the inte-
gration table. According to Dormand [24], the Gauss-Jackson formulae have been
developed for the convenience of manual operation and offer little or no advantage
over the Störmer-Cowell formulae, which lend themselves to be easily programmed
on a computer. Without objecting to this, we observe that the Gaussian integration
has been used through two centuries by the astronomers well before the advent of
modern computing means. In addition, according to Henrici [2] and Herrick [55,
56], the Gauss-Jackson method is less sensitive than the Störmer-Cowell to
roundoff errors resulting from a computation carried out with limited number of
decimal places.
An example of application of the Gauss-Jackson method is given below. Let us
consider the following differential equation
x00 ¼ x
defined for all values of the independent variable t with the following initial con-
ditions x0 = 0, x′0 = 1. As is well known, the analytical solution of this initial-value
problem is x(t) = sin t. We want to integrate numerically this differential equation,

with the specified initial conditions, by means of a Gauss-Jackson fourth-order
formula, with a step size h = 0.1. Let the following starting values be known
x0:3 ¼ 0:2955202
x0:2 ¼ 0:1986693
x0:1 ¼ 0:0998334
x0:0 ¼ 0:0000000
x0:1 ¼ 0:0998334
x0:2 ¼ 0:1986693
x0:3 ¼ 0:2955202
where the subscripts are the corresponding values of t. To integrate numerically the
given differential equation, we use the following integration table
i ti ∇−2fi ∇−1fi fi ∇fi ∇2fi ∇3fi ∇4fi

−3 −0.3
−2 −0.2
−1 −0.1
0 0.0
1 0.1
2 0.2
3 0.3
Since in the present case the integrand function f(x) is such that fi = −xi, then we
can fill the column of fi in the table by writing there the same values (but with the
minus sign) as those contained in the column of xi, as shown below.

−3 −0.3 0.2955202
−2 −0.2 0.1986693
−1 −0.1 0.0998334
0 0.0 0.0000000
1 0.1 −0.0998334
2 0.2 −0.1986693
3 0.3 −0.2955202
Now, operating as has been shown in Sects. 6.14 and 6.15, we can fill the columns
of the first, second, third, and fourth differences, as shown in the following table.

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
Now, we use the formula given above, that is,

2 xi 1 1 2 1 3 221 4
r fi1 ¼ 2 r r r fi
h 12 240 240 60480
with, sequentially, i = 3, i = 2, and i = 1, in order to compute ∇−2f2, ∇−2f1 and

∇−2f0, as shown below.
x3 1 1 2 1 3 221 4 0:2955202
r2 f2 ¼ f3 þ r f3 þ r f3 þ r f3 ¼
h2 12 240 240 60480 0:12
0:2955202 0:0019850 0:0009875 221
þ þ þ 0:0000100
12 240 240 60480
¼ 29:5766598
x2 1 1 2 1 3 221 0:1986693
r2 f1 ¼ f2 þ r f2 þ r f2 þ r4 f2 ¼
h2 12 240 240 60;480 0:12
0:1986693 0:0009975 0:0009975 221
þ þ þ 0 ¼ 19:8834972
12 240 240 60480
x1 1 1 2 1 3 221 0:0998334
r2 f0 ¼ 2 f1 þ r f1 þ r f1 þ r4 f1 ¼
h 12 240 240 60;480 0:12
0:0998334 0 0:0009975 221
þ þ þ 0:00001 ¼ 9:9916653
12 240 240 60;480
These values are written in the column of ∇−2fi in the integration table, as
follows.

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
The values of ∇−2fi are used to compute the correspondent values of ∇−1fi,
according to the expression given above ∇−1fi = ∇−2fi − ∇−2fi−1, for i = 1, 2.
Thus, ∇−1f1 and ∇−1f2 result from
r1 f1 ¼ r2 f1 r2 f0 ¼ 19:8834972 9:9916653 ¼ 9:8918319

r1 f2 ¼ r2 f2 r2 f1 ¼ 29:5766598 19:8834972 ¼ 9:6931626
These values are written in the column of ∇−1fi in the table given above. The
results are shown in the following table.
i ti ∇−2fi ∇−1fi fi ∇fi ∇2fi ∇3 f i ∇4fi

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 9.9818319 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 9.6931626 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
Since f3 and ∇−1f2 are known, then ∇−1f3 can be computed as follows
r1 f3 ¼ f3 þ r1 f2 ¼ 0:2955202 þ 9:6931626 ¼ 9:3976424
This is because, as shown above, fi = ∇−1fi − ∇−1fi−1.

Now, since ∇−1f3 and ∇−2f2 are known, then ∇−2f3 can be computed as follows
r2 f3 ¼ r1 f3 þ r2 f2 ¼ 9:3976424 þ 29:5766598 ¼ 38:9743022
By inserting the values found above in the table, we have
i ti ∇−2fi ∇−1fi fi ∇fi ∇2fi ∇3 f i ∇4fi

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 9.9818319 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 9.6931626 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 38.9743022 9.3976424 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
We now apply the Gauss-Jackson predictor formula given above, which is

ð pÞ 2 1 1 19 2 3 3 863 4
xi þ 1 ¼h 2
r fi þ fi þ rfi þ r fi þ r fi þ r fi
12 12 240 40 12096
Thus, we have

ð pÞ 1 1 19 2 3 3 863 4
x4 ¼ h2 r2 f3 þ f3 þ rf3 þ r f3 þ r f3 þ r f3
12 12 240 40 12096
0:2955202 0:0968509 19
¼ 0:12 ð38:9743022 þ 0:0019850
12 12 240
3 863
þ 0:0009875 0:00001Þ ¼ 0:3894184
40 12096
We can now evaluate the integrand function at t = 0.4 by means of the value of x(p)
4
resulting from the predictor formula, as follows

ðpÞ ðpÞ
f4 ¼ f t4 ; x4 ¼ x4 ¼ 0:3894184
By inserting this value in the table, we have

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 9.9818319 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 9.6931626 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 38.9743022 9.3976424 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
4 0.4 −0.3894184
Now, we can compute the differences of f4 as follows
rf4 ¼ f4 f3 ¼ 0:3894184 ð0:2955202Þ ¼ 0:0938982

r2 f4 ¼ rf4 rf3 ¼ 0:0938982 ð0:0968509Þ ¼ 0:0029527
r3 f4 ¼ r2 f4 r2 f3 ¼ 0:0029527 0:0019850 ¼ 0:0009677
r4 f4 ¼ r3 f4 r3 f3 ¼ 0:0009677 0:0009875 ¼ 0:0000198
By inserting these values in the table, we have

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 9.9818319 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 9.6931626 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 38.9743022 9.3976424 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
4 0.4 −0.3894184 −0.0938982 0.0029527 0.0009677 −0.0000198
Now, since f4 and ∇−1f3 are known, ∇−1f4 can be computed as follows
r1 f4 ¼ f4 þ r1 f3 ¼ 0:3894184 þ 9:3976424 ¼ 9:0082240
Likewise, since ∇−1f4 and ∇−2f3 are known, ∇−2f4 can be computed as follows
r2 f4 ¼ r1 f4 þ r2 f3 ¼ 9:0082240 þ 38:9743022 ¼ 47:9825262
By inserting these values in the table, we have

−3 −0.3 0.2955202
−2 −0.2 0.1986693 −0.0968509
−1 −0.1 0.0998334 −0.0988359 −0.0019850
0 0.0 9.9916653 0.0000000 −0.0998334 −0.0009975 0.0009875
1 0.1 19.8834972 9.9818319 −0.0998334 −0.0998334 0.0000000 0.0009975 0.0000100
2 0.2 29.5766598 9.6931626 −0.1986693 −0.0988359 0.0009975 0.0009975 0.0000000
3 0.3 38.9743022 9.3976424 −0.2955202 −0.0968509 0.0019850 0.0009875 −0.0000100
4 0.4 47.9825262 9.0082240 −0.3894184 −0.0938982 0.0029527 0.0009677 −0.0000198
Now, x4 can be computed again by means of the Gauss-Jackson corrector

formula

ðcÞ 1 1 2 1 3 221 4
xi ¼ h2 r2 fi1 þ fi r fi r fi r fi
12 240 240 60480
as follows

ð cÞ 2 1 1 2 1 3 221 4
x4 ¼ h r f3 þ
2
f4 r f4 r f4 r f4
12 240 240 60480
0:3894184 0:0029527 0:0009677
¼ 0:12 ð38:9743022
12 240 240
221
0:0000198Þ ¼ 0:3894183
60480
The value of x4 computed analytically, with seven decimal places, is
x4 ¼ sinð0:4Þ ¼ 0:3894183
The value of x′4 can be computed by means of the Adams-Moulton corrector

formula of Sect. 6.14, which, truncated after the fourth difference, is

0ðcÞ 1 1 1 19 3 3 4
xi ¼ h r1 fi fi rfi r2 fi r fi r fi
2 12 24 720 160
In the example considered above, there results

0 ðcÞ 1 1 1 19 3 3 4
x4 ¼ h r1 f4 f4 rf4 r2 f4 r f4 r f4
2 12 24 720 160

0:3894184 0:0938982 0:0029527
¼ 0:1 9:0082240 þ þ
2 12 24

19 3
0:0009677 þ 0:0000198 ¼ 0:9210610
720 160
The value of x′4 computed analytically, with seven decimal places, is
x0 4 ¼ cosð0:4Þ ¼ 0:9210610
6.17 Calculation of the Starting Values
As has been shown in Sect. 6.13, the multi-step methods are not self-starting. In
order to start any one of such methods, a single-step method is generally required
(see Sects. 6.2–6.12). The present paragraph describes some techniques for com-
puting the required starting values.
(a) Single-step methods
These methods do not require a set of initial points to start. Bearing in mind the
necessity of using a single-step method having the same accuracy as that of the
multi-step method to be used afterwards, the Runge–Kutta method is usually
employed. For example, let us use the Runge–Kutta method to compute the four
additional starting values which are necessary to integrate numerically the fol-
lowing differential equation (from [51])
x0 ¼ t þ x
with the initial condition t0 = 0, x0 = 1, with a step size h = 0.05.

6.17 Calculation of the Starting Values 983
The following Runge–Kutta formulae of Sect. 6.2 will be used below:

1 1
2 2

1 1
2 2
k4 ¼ f ðtn þ h; xn þ hk3 Þ
1
U ¼ ðk1 þ 2k2 þ 2k3 þ k4 Þ
6
1
6
(1) First step forward
t¼0
x¼1
k1 ¼ f ðt; xÞ ¼ 0 þ 1 ¼ 1
t ¼ 0 þ 0:05=2 ¼ 0:025
x ¼ 1 þ 0:05=2 ¼ 1:025
k2 ¼ f ðt; xÞ ¼ 0:025 þ 1:025 ¼ 1:05
t ¼ 0:025
x ¼ 1 þ 0:05 1:05=2 ¼ 1:02625
k3 ¼ f ðt; xÞ ¼ 0:025 þ 1:02625 ¼ 1:05125
t ¼ 0 þ 0:05 ¼ 0:05
x ¼ 1 þ 0:05 1:05125 ¼ 1:0525625
k4 ¼ f ðt; xÞ ¼ 0:05 þ 1:0525625 ¼ 1:1025625
U ¼ ð1 þ 2 1:05 þ 2 1:05125 þ 1:1025625Þ=6 ¼ 1:05084375
x1 ¼ 1 þ 0:05 1:05084375 ¼ 1:052542188
(2) Second step forward
t ¼ 0:05
x ¼ 1:052542188
k1 ¼ f ðt; xÞ ¼ 0:05 þ 1:052542188 ¼ 1:102542188
t ¼ 0:05 þ 0:05=2 ¼ 0:075
x ¼ 1:052542188 þ 0:05 1:102542188=2 ¼ 1:080105743
k2 ¼ f ðt; xÞ ¼ 0:075 þ 1:080105743 ¼ 1:155105743
t ¼ 0:075
x ¼ 1:052542188 þ 0:05 1:155105743=2 ¼ 1:081419832
k3 ¼ f ðt; xÞ ¼ 0:075 þ 1:081419832 ¼ 1:156419832
t ¼ 0:05 þ 0:05 ¼ 0:1
x ¼ 1:052542188 þ 0:05 1:156419832 ¼ 1:11036318
k4 ¼ f ðt; xÞ ¼ 0:1 þ 1:11036318 ¼ 1:21036318
U ¼ ð1:102542188 þ 2 1:155105743 þ 2 1:156419832 þ 1:21036318Þ=6
¼ 1:155992753
x2 ¼ x1 þ h U ¼ 1:052542188 þ 0:05 1:155992753 ¼ 1:110341826
(3) First step backward
t¼0
x¼1
k1 ¼ f ðt; xÞ ¼ 0 þ 1 ¼ 1
t ¼ 0 0:05=2 ¼ 0:025
x ¼ 1 0:05=2 ¼ 0:975
k2 ¼ f ðt; xÞ ¼ 0:025 þ 0:975 ¼ 0:95
t ¼ 0:025
x ¼ 1 0:05 0:95=2 ¼ 0:97625
k3 ¼ f ðt; xÞ ¼ 0:025 þ 0:97625 ¼ 0:95125
t ¼ 0 0:05 ¼ 0:05
x ¼ 1 0:05 0:95125 ¼ 0:9524375
k4 ¼ f ðt; xÞ ¼ 0:05 þ 0:9524375 ¼ 0:9024375
U ¼ ð1 þ 2 0:95 þ 2 095125 þ 0:9024375Þ=6 ¼ 0:950822917
x1 ¼ 1 0:05 0:950822917 ¼ 0:952458854
(4) Second step backward
t ¼ 0:05
x ¼ 0:952458854
k1 ¼ f ðt; xÞ ¼ 0:05 þ 0:952458854 ¼ 0:902458854
t ¼ 0:05 0:05=2 ¼ 0:075
x ¼ 0:952458854 0:05 0:902458854=2 ¼ 0:929897383
k2 ¼ f ðt; xÞ ¼ 0:075 þ 0:929897383 ¼ 0:854897383
t ¼ 0:075
x ¼ 0:952458854 0:05 0:854897383=2 ¼ 0:931086419
k3 ¼ f ðt; xÞ ¼ 0:075 þ 0:931086419 ¼ 0:856086419
t ¼ 0:05 0:05 ¼ 0:1
x ¼ 0:952458854 0:05 0:856086419 ¼ 0:909654533
k4 ¼ f ðt; xÞ ¼ 0:1 þ 0:909654533 ¼ 0:809654533
U ¼ ð0:902458854 þ 2 0:854897383 þ 2 0:856086419 þ 0:809654533Þ=6
¼ 0:855680165
x2 ¼ x1 þ h U ¼ 0:952458854 0:05 0:855680165 ¼ 0:909674846
This accounts for the starting values given in Sect. 6.14.

(b) Taylor-series expansions
If the derivatives f′, f″, … of the integrand function f(t, x) can easily be com-
puted, then the necessary starting values can be determined by expanding the
unknown function x(t) in a Taylor series (see [5]):
1 00 1
xðnhÞ xðt0 Þ þ x0 ðt0 Þnh þ x ðt0 ÞðnhÞ2 þ xð3Þ ðt0 ÞðnhÞ3 þ
2! 3!
(where n = ± 1, ± 2, …), using as many terms on the right-hand side before

truncation as are required.
An example (taken from [51]) shows how to start the solution, approximated to
the fourth order, of the differential equation x′ = t + x, with the initial condition x
(0) = x0 = 1. First, we compute the derivatives of x, as follows
x0 ¼ t þ x x0 0 ¼ t0 þ x0 ¼ 1
x00 ¼ 1 þ x0 x00 0 ¼ 1 þ x0 0 ¼ 1 þ 1 ¼ 2
x000 ¼ x00 x000 0 ¼ x00 0 ¼ 2
xIV ¼ x00 xIV 0 ¼ x000 0 ¼ 2
xV ¼ xIV xV 0 ¼ xIV 0 ¼ 2
Then, we expand x(t) in a Taylor series, truncated after the fifth power of h, with
h = 0.05 and n = ± 1, n = ± 2, as follows
1 00 2 1 ð3Þ 3 1 ð4Þ 4 1 ð5Þ 5

x 1 ¼ x ð t 0 þ hÞ ¼ x 0 þ x 0 0 h þ
x 0 h þ x 0 h þ x0 h þ x0 h
2! 3! 4! 5!
0:052 0:053 0:054 0:055
¼ 1 þ 1 0:05 þ 2 þ2 þ2 þ2
2 32 432 5432
¼ 1:0525421927
1 00 2 1 ð3Þ 3 1 ð4Þ 4 1 ð5Þ 5

x1 ¼ xðt0 hÞ ¼ x0 x0 0 h þ x 0 h x0 h þ x0 h x0 h
2! 3! 4! 5!
0:052 0:053 0:054 0:055
¼ 1 1 0:05 þ 2 2 þ2 2
2 32 432 5432
¼ 0:952458848958
1 00 1 ð3Þ 1 ð4Þ
x2 ¼ xðt0 þ 2hÞ ¼ x0 þ x0 0 ð2hÞ þ
x 0 ð2hÞ2 þ x0 ð2hÞ3 þ x0 ð2hÞ4
2! 3! 4!
1 ð5Þ 0:12 0:13 0:14
þ x0 ð2hÞ5 ¼ 1 þ 1 0:1 þ 2 þ2 þ2
5! 2 32 432
0:15
þ 2 ¼ 1:1103418333
5432
1 0 1 ð3Þ 1 ð4Þ
x2 ¼ xðt0 2hÞ ¼ x0 x0 0 ð2hÞ þ
x 0 ð2hÞ2 x0 ð2hÞ3 þ x0 ð2hÞ4
2! 3! 4!
1 ð5Þ 0:12 0:13 0:14
x0 ð2hÞ5 ¼ 1 1 0:1 þ 2 2 þ2
5! 2 32 432
0:15
2 ¼ 0:9096748333
5432
Substituting these values in the given differential equation, we have the following
five consecutive values of x′ needed to start the solution
x0 2 ¼ t2 þ x2 ¼ 0:1 þ 0:9096748333 ¼ 0:809674833

x0 1 ¼ t1 þ x1 ¼ 0:05 þ 0:952458848958 ¼ 0:902458848958
x0 0 ¼ t0 þ x0 ¼ 0 þ 1 ¼ 1
x0 1 ¼ t1 þ x1 ¼ 0:05 þ 1:0525421927 ¼ 1:1025421927
x0 2 ¼ t2 þ x2 ¼ 0:1 þ 1:1103418333 ¼ 1:2103418333
Should the function x(t) be not-existent for values of t less than t0, then the com-
putation would be performed only for t greater than t0, by substituting the proper
values of t and x in the given differential equation.
(c) Interpolating formulae

Another starting procedure for a fourth-order method is suggested by Ralston
[57]. This procedure is based on Newton’s interpolation formula and uses forward
differences to get the required values of x1, x2, and x3 as a linear combination of x0
and x′0, x′1, x′2, and x′3, that is,
1
x1 ¼ x0 þ hð9x0 0 þ 19x0 1 5x0 2 þ x0 3 Þ
24
1
x2 ¼ x0 þ hðx0 0 þ 4x0 1 þ x0 2 Þ
3
3
x3 ¼ x0 þ hðx0 0 þ 3x0 1 þ 3x0 2 þ x0 3 Þ
8
The error in the three equations written above is O(h5). To use them, we first
estimate x′1, x′2, and x′3 by means of the simple Euler method and then compute x1,
x2, and x3. These values are introduced in the differential equation to compute
improved values of x′1, x′2, and x′3, which in turn are used to compute again x1, x2,
and x3, until convergence is reached. To show an application of this starting pro-
cedure, let us consider the following differential equation
tx
x0 ¼
t 2 þ x2
with the initial condition x(0) = 1. First, we compute x″, as follows
ðx þ t x0 Þðt2 þ x2 Þ t xð2t þ 2x x0 Þ
x00 ¼
ð t 2 þ x2 Þ 2
when t = 0, we have x′0 = 0 and x″0 = 1.

We take h = 0.1 and apply Euler’s method to compute trial values of x′1, x′2, and
x′3, as follows
x0 1 ¼ x0 0 þ hx00 0 ¼ 0 þ 0:1 1 ¼ 0:1

x0 2 ¼ x0 0 þ 2hx00 0 ¼ 0 þ 2 0:1 1 ¼ 0:2
x0 3 ¼ x0 0 þ 3hx00 0 ¼ 0 þ 3 0:1 1 ¼ 0:3
Now, we use x0 and the trial values of x′1, x′2, and x′3 to compute first-iteration
values of x1, x2, and x3, as follows

9 0 þ 19 0:1000 5 0:2000 þ 0:3000
x1 ¼ 1:0000 þ 0:1000 ¼ 1:0050
24

0 þ 4 0:1000 þ 0:2000
x2 ¼ 1:0000 þ 0:1000 ¼ 1:0200
3
3
x3 ¼ 1:0000 þ 0:1000 ð0 þ 3 0:1000 þ 3 0:2000 þ 0:3000Þ ¼ 1:0450
8
These values are substituted into the differential equation to compute improved
values of x′1, x′2, and x′3, as follows
0:1000 1:0050
x0 1 ¼ ¼ 0:0985
0:10002 þ 1:00502
0:2000 1:0200
x0 2 ¼ ¼ 0:1888
0:20002 þ 1:02002
0:3000 1:0450
x0 3 ¼ ¼ 0:2652
0:30002 þ 1:04502
Now, the values of x1, x2, and x3 are computed again, as follows

9 0 þ 19 0:0985 5 0:1888 þ 0:2652
x1 ¼ 1:0000 þ 0:1000 ¼ 1:0050
24

0 þ 4 0:0985 þ 0:1888
x2 ¼ 1:0000 þ 0:1000 ¼ 1:0194
3
3
x3 ¼ 1:0000 þ 0:1000 ð0 þ 3 0:0985 þ 3 0:1888 þ 0:2652Þ ¼ 1:0423
8
Now, these values are substituted again into the differential equation to compute
improved values of x′1, x′2, and x′3, as follows
0:1000 1:0050
x0 1 ¼ ¼ 0:0985
0:10002 þ 1:00502
0:2000 1:0194
x0 2 ¼ ¼ 0:1889
0:20002 þ 1:01942
0:3000 1:0423
x0 3 ¼ ¼ 0:2658
0:30002 þ 1:04232
Now the values of x1, x2, and x3 are computed again, as follows

9 0 þ 19 0:0985 5 0:1889 þ 0:2658
x1 ¼ 1:0000 þ 0:1000 ¼ 1:0050
24

0 þ 4 0:0985 þ 0:1889
x2 ¼ 1:0000 þ 0:1000 ¼ 1:0194
3
3
x3 ¼ 1:0000 þ 0:1000 ð0 þ 3 0:0985 þ 3 0:1889 þ 0:2658Þ ¼ 1:0423
8
Since the values of x1, x2, and x3 computed here are equal, within the desired
tolerance, to those computed in the previous iteration, they are taken as correct.
(d) Starting procedures for the special case of the n-body problem
This technique consists in taking an initial guess (or prediction) of the
back-points and refining the values coming from this guess by means of corrector
formulae. Such formulae, which obtain improved values by using not only previous
points but also subsequent points, are called mid-corrector formulae.
The guess which predicts the initial values is provided by the analytical solution
of the corresponding two-body problem. This technique is fully described in [52].
6.18 Halving the Step Size
The operation of reducing the step size is the same as that of finding starting values.
To do this by using fourth-order formulae, we interpolate five consecutive values of
the dependent variable vector x previously computed and get two more ordinates by
means of the fourth-degree interpolating polynomial. This polynomial is unique,
independently of the method used for determining it, because a fourth-degree
polynomial x = at4 + bt3 + ct2 + dt + e is determined by the values of its five
coefficients a, b, c, d, and e. By imposing to this polynomial to pass through five
points of given co-ordinates x and t, we have a set of five equations for the five
unknowns a, b, c, d, and e. Generally speaking, in the case of a k-step method,
k + 1 previous points are required to restart the method by using a halved inte-
gration step size. In the case of a fourth-order method, we can use the
Lagrangian interpolation through the points (ti, xi), (ti−1, xi−1), (ti−2, xi−2), (ti−3,
xi−3), and (ti−4, xi−4), equispaced by h.
By so doing, it possible to compute two intermediate points (ti−1/2, xi−1/2) and
(ti−3/2, xi−3/2). The formulae (from [46]) to be used with a fourth-order method are
given below:
1
xi12 ¼ ð35xi þ 140xi1 70xi2 þ 28xi3 5xi4 Þ
128
1
xi32 ¼ ðxi þ 24xi1 þ 54xi2 16xi3 þ 3xi4 Þ
64
We have thus a further set of five equidistant ordinates (xi, xi−1/2, xi−1, xi−3/2 and
xi−2), equispaced by Dt = ½h. Now, it is possible to start again the integration with
a new step size equal to a half of the previous step size.
Formulae for Lagrangian interpolation with a variety of degrees may be found in
[58]. This method of interpolating the ordinates lends itself to be easily programmed
on a computer. If a numerical integration is performed by means of a table of
differences, another way for finding new values at the midpoints of intervals pre-
viously computed is based on the use of Bessel’s formula for interpolating to halves:

0 x0 0 þ x01 1 r2 x0 1 þ r2 x0 2 3 r 4 x0 2 þ r 4 x0 3
x 12 ¼ þ
2 8 2 128 2
6 0 6 0

5 r x 3þr x 4 ½1 3 5 ð2n 1Þ2
þ þ ð1Þn
1024 2 22n ð2nÞ!
where x′½ is the midpoint between x′0 and x′1 and the differences ∇2, ∇4, ∇6, …, ∇2n
are ordinary backward differences.
In the case of a fourth-degree polynomial, the preceding formula reduces to

x0 0 þ x0 1 1 r 2 x0 1 þ r 2 x0 2 3 r 4 x0 2 þ r 4 x0 3
x0 12 ¼ þ
2 8 2 128 2
and it is necessary to compute the value of x′ at two mid-intervals in order to have

five consecutive values of x′ to be used for restarting the computation with halved
step size. An example, taken from [51], illustrates this method.
Let us consider again the differential equation
x0 ¼ t þ x
with the initial condition t0 = 0, x0 = 1. Suppose that, having computed the fol-
lowing table
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

0.0 1.0000 1.0000
0.1 1.1103 1.2103 0.2103
0.2 1.2427 1.4427 0.2324 0.0221
0.3 1.3996 1.6996 0.2569 0.0245 0.0024
0.4 1.5835 1.9835 0.2839 0.0270 0.0025 0.0001
0.5 1.7973 2.2973 0.3138 0.0299 0.0029 0.0004
0 0.6 2.0441 2.6441 0.3468 0.0330 0.0031 0.0002
1 0.7 2.3274 3.0274 0.3833 0.0365 0.0035 0.0004
2 0.8 2.6510 3.4510 0.4236 0.0403 0.0038 0.0003
3 0.9 3.0191 3.9191 0.4681 0.0445 0.0042 0.0004
4 1.0 3.4365 4.4365 0.5174 0.0493 0.0048 0.0006
we want to halve the step size from t = 0.8 onwards.

To this end, we apply Bessel’s formula truncated after the fourth difference,
starting from t = 0.6 (i.e., t = 0.6 is the zero point), as follows

2:6441 þ 3:0274 1 0:0365 þ 0:0403 3 0:0003 þ 0:0004
x0 0:65 ¼ þ
2 8 2 128 2
¼ 2:8310
6.18 Halving the Step Size 991
Then, we apply the same formula, starting from t = 0.7 (i.e., t = 0.7 is the zero
point), as follows

3:0274 þ 3:4510 1 0:0403 þ 0:0445 3 0:0004 þ 0:0006
x0 0:75 ¼ þ
2 8 2 128 2
¼ 3:2339
The five starting values for the new table are three (x′0.60, x′0.70, and x′0.80) of those
computed previously by numerical integration with step size h = 0.10 plus the two
(x′0.65 and x′0.75) computed by interpolation.
They fill the column of x′ in the new table, as shown below.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

0.60 2.6441
0.65 2.8310
0.70 3.0274
0.75 3.2339
0.80 2.6510 3.4510
Now, we can compute ∇x′, ∇2x′, ∇3x′, and ∇4x′, as shown in Sect. 6.14, and the
results are used to fill the corresponding columns of the new table, as shown below.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

0.60 2.6441
0.65 2.8310 0.1869
0.70 3.0274 0.1964 0.0095
0.75 3.2339 0.2065 0.0101 0.0006
0.80 2.6510 3.4510 0.2171 0.0106 0.0005 −0.0001
The computation can be started again from t = 0.8 with the step size halved.
For this purpose, we apply the Adams predictor formula to the line corre-
sponding to t = 0.80 and compute xp0.85; this predicted value is used to compute the
derivative x′0.85. The results are given below.

0 1 0 5 2 0 3 3 0 251 4 0
xp0:85 ¼ x0:80 þ h x 0:80 þ rx 0:80 þ r x 0:80 þ r x 0:80 þ r x 0:80
2 12 8 720
After substituting h = 0.05 and the values indicated in the last row
(x0.80 = 2.6510, x′0.80 = 3.4510, ∇x′0.80 = 0.2171, ∇2x′0.80 = 0.0106, ∇3x′0.80 =
0.0005, ∇4x′0.80 = −0.0001) of the preceding table, we find
xp0:85 ¼ 2:8292
The corresponding value of the first derivative is

x0 0:85 ¼ f 0:85; xp0:85 ¼ 0:85 þ xp0:85 ¼ 0:85 þ 2:8292 ¼ 3:6792
We use the value of x′0.85 computed above to fill the line of the table which
corresponds to t = 0.85, as shown below.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

0.80 2.6510 3.4510 0.2171 0.0106 0.0005 −0.0001
0.85 3.6792 0.2282 0.0111 0.0005 0.0000
0.90
0.95
1.00
and apply the corrector as a check, as follows

1 1 1 19 4 0
xc0:85 ¼ x0:80 þ h x 0:85 rx0 0:85 r2 x0 0:85 r3 x0 0:85
0
r x 0:85
2 12 24 720
By substituting h = 0.05, x0.80 = 2.6510, and the values indicated in the third
row (x′0.85 = 3.6792, ∇x′0.85 = 0.2282, ∇2x′0.85 = 0.0111, ∇3x′0.85 = 0.0005, and
∇4x′0.85 = 0) of the preceding table, we find
xciþ 1 ¼ 2:8292
This value is the same as that computed previously by means of the predictor
and can be used to fill the row corresponding to t = 0.85, as shown below.
t x x′ ∇x′ ∇2x′ ∇3x′ ∇4x′

0.80 2.6510 3.4510 0.2171 0.0106 0.0005 −0.0001
0.85 2.8292 3.6792 0.2282 0.0111 0.0005 0.0000
0.90
0.95
1.00
Now, we apply the predictor to the line corresponding to t = 0.85, as follows

1 5 2 0 3 251 4 0
xp0:90 ¼ x0:85 þ h x0 0:85 þ rx0 0:85 þ r x 0:85 þ r3 x0 0:85 þ r x 0:85
2 12 8 720
6.18 Halving the Step Size 993
By substituting h = 0.05 and the values indicated in the row corresponding

to t = 0.85 (x0.85 = 2.8292, x′0.85 = 3.6792, ∇x′0.85 = 0.2282, ∇2x′0.85 = 0.0111,
∇3x′0.85 = 0.0005, and ∇4x′0.85 = 0) of the preceding table, we find
xp0:90 ¼ 3:0191
and compute the derivative x′0.90 by means of this value, as follows

x0 0:90 ¼ f 0:90; xp0:90 ¼ 0:90 þ xp0:90 ¼ 0:90 þ 3:0191 ¼ 3:9191
These are the same values as those computed before halving the step size.
Therefore, no error has been induced by changing to a smaller step size. The shorter
interval reduces the fourth differences to negligible values.
Another procedure which can be used to halve the step size has been suggested
by Ralston [57]. This procedure is based on the following formulae
1
xi12 ¼ ½80xi þ 135xi1 þ 40xi2 þ xi3 þ hð15x0 i þ 90x0 i1 þ 15x0 i2 Þ
256
1
xi32 ¼ ½12xi þ 135xi1 þ 108xi2 þ xi3 þ hð3x0 i 5x0 i1 þ 27x0 i2 Þ
256
where the error term in the two expressions given above is O(h7).
As an example of application, we consider once again the differential equation of
Sect. 6.14 (x′ = t + x, with the initial condition t0 = 0, x0 = 1). We want to halve
the step size h = 0.1 from t = 0.8 onwards. For this purpose, we set i = 2 and use
the tables shown above to get
xi ¼ 2:6510 x0 i ¼ 3:4510
xi1 ¼ 2:3274 x0 i1 ¼ 3:0274
xi2 ¼ 2:0441 x0 i2 ¼ 2:6441
xi3 ¼ 1:7973
Then, we apply the formulae indicated above as follows
1
x0:75 ¼ ½80 2:6510 þ 135 2:3274 þ 40 2:0441 þ 1:7973
256
þ 0:1 ð15 3:4510 þ 90 3:0274 þ 15 2:6441Þ ¼ 2:4839
1
x0:65 ¼ ½12 2:6510 þ 135 2:3274 þ 108 2:0441 þ 1:7973
256
þ 0:1 ð3 3:4510 54 3:0274 þ 27 2:6441Þ ¼ 2:1810
By inspection of the difference tables, these values turn out to be the same as
those computed previously.
Of course, a further way for decreasing the step size is to compute new points by
means of a single-step method (e.g., by means of the Runge–Kutta method).
6.19 Integration for Elliptic Orbits of High Eccentricity
For this type of orbits, integration with fixed step size is less efficient than that with
variable step size. This is because a small step length is required when the orbiting
object is near the centre of attraction, and therefore moves rapidly, whereas a large
step length is sufficient when the same object is far away from the attracting body,
and therefore moves slowly.
If the integration is performed by means of the Gauss-Jackson method (described
in Sect. 6.16), Merson [54] has suggested that time (t) should be transformed to a
new independent variable (s) given by
dt 3
¼ r2
ds
By so doing, constant changes in s generate small changes in t when the radius

vector (r) is small, and large changes in t when r is large. The integration can be
performed by using this method (which Fox in [59] calls Gauss–Jackson–Merson
method) with a constant step size Ds over the whole length of the interval. This
method has the following disadvantages:
• the necessity of solving the additional equation dt/ds = r3/2, to transform back
from s to t, for the purpose of computing the time at each step;
• the solution being required up to some given time, it is unlikely that the equi-
spaced intervals Ds should produce a set of results just at this desired time, with
consequent necessity of interpolating the results at the endpoint; and
• the necessity of using the corrector mode of the Gauss–Jackson method to
ensure the stability of the solution.
However, it is to be observed that:
• the work necessary to solve the additional equation dt/ds = r3/2 is negligible in
comparison with that required by the equations of motion;
• the interpolation can be performed easily, because the required difference tables
have been computed during the integration process; and
• the necessity of applying the corrector produces a small increase in the time
required to advance from a step to the next.
For example, the Gauss–Jackson predictor formulae (see Sect. 6.16) for an
eighth-order integration can be written as follows
6.19 Integration for Elliptic Orbits of High Eccentricity 995

1 00 1 19 2 00 3 3 00
xpiþ 1 xðti þ hÞ ¼ h r2 x00 i þ
2
x iþ rx00 i þ r x iþ r x i
12 12 240 40

863 4 00 275 5 00 33953 6 00 8183 7 00
þ r x iþ r x iþ r x iþ r x i
12096 4032 518400 129600

1 5 3 251 3 00
yi þ 1 yðti þ hÞ ¼ h r1 x00 i þ x00 i þ
p
rx00 i þ r2 x00 i þ r x i
2 12 8 720

95 4 00 19087 5 00 5257 6 00 1070017 7 00
þ r x iþ r x iþ r x iþ r x i
288 60480 17280 3628800
in the backward difference notation, where y x′. If the accelerations do not

depend on the velocities, that is, if x′′ = f(t, x), then the velocity vector y need not
be computed at each step. Following Fox [59] and Herrick [55, 56], the interpo-
lation formulae are

2 00 1 00 1 1 00 1 3 r
xðti þ r hÞ ¼ h r x i1 þ rr x i12 þ
2
r þ
2
x iþ r
2! 6 3! 2

1 1 1 10 11
rx00 i þ 12 þ r4 r2 x00 i þ 1 þ r5 r3 þ r
4! 10 5! 3 6

1 5 31 1 21
r3 x00 i þ 32 þ r6 r4 þ r4 x00 i þ 2 þ r7 r5
6! 2 84 7! 2

191 1 28 53 4 289
þ 28r3 r r5 xi þ 52 þ r8 r6 þ r
12 8! 3 3 90

1 882 5 2497
r6 x00 i þ 3 þ r9 24 r7 þ r 432 r3 þ r r4 x00 i þ 72
9! 5 10

1 1 1
yðti þ r hÞ ¼ h r1 x00 i12 þ rx00 i þ r2 rx00 i þ 12 þ r3 r2 x00 i þ 1
2! 6 3!

1 11 1 5
þ r4 2 r2 þ r3 x00 i þ 32 þ r5 r3 r4 x00 i þ 2
4! 30 5! 3

1 15 4 191 5 00 1 28 3
þ r r þ 12 r
6 2
r x i þ 52 þ r7 7r5 þ r r6 x00 i þ 3
6! 2 84 7! 3

1 56 6 2497 7 00
þ r r þ 98 r 144 r þ
8 4 2
r x i þ 72
8! 3 90
where r = (t − ti)/(ti+1 − ti) is the interpolating factor, and x′′i+½ = ½(x′′i+1 + x′′i),
….
It is to be noted that the integration must be extended to ti+4 before the inter-
polation in the range [ti, ti+1] can be done.
The correspondent Gauss-Jackson corrector formulae (see Sect. 6.16) are

1 00 1 2 00 1 3 00
xi þ 1 xðti þ hÞ ¼ h r2 x00 i þ
c 2
x iþ1 r x iþ1 r x iþ1
12 240 240

221 4 00 19 5 00 9829 407
r x iþ1 r x iþ1 r6 x00 i þ 1 r7 x00 i þ 1
60480 6048 3628800 172800

1 1 1
yciþ 1 yðti þ hÞ ¼ h r1 x00 i þ x00 i þ 1 rx00 i þ 1 r2 x00 i þ 1
2 12 24

19 3 00 3 4 00 863 5 00 275 6 00 33953 7 00
r x iþ1 r x iþ1 r x iþ1 r x iþ1 r xi þ 1
720 160 60480 24192 3628800
Sundman [60] and successively Nacozy [61] and Berry et al. [62] have modified the
time transformation dt/ds = r3/2 by choosing
dt
¼ cr n
ds
which is called the generalised Sundman time transformation, where

n ¼ 3=2
c ¼ l1=2
s is called the intermediate anomaly of Keplerian motion
l = GM is the gravitational parameter of the Earth
G is the gravitational constant
M is the mass of the Earth.
This procedure is called the s-integration, which can be summarised as follows.
First, a step size is selected in terms of time. The time step size (Dt) must be
converted to the intermediate anomaly step size (Ds) by means of
1 3
Ds ¼ l2 rp 2 Dt
where rp is the distance between the satellite at perigee and the Earth centre. Having
converted the step size in terms of intermediate anomaly, the integration process
enters its start-up phase, by computing eight positions and velocities around epoch,
by using f and g series. These series depend on the time between the points, and the
points must be equispaced in s, in order for the s-integration to be performed.
Therefore, it is necessary to convert from intermediate anomaly to time. For this
purpose, in this initial phase, the time is approximated (see [52]) by holding the
epoch distance constant, as follows
3
Dtfg ¼ r02 l2 Ds
1
Now, the eight position and velocities are computed by means of the f and g series,
and this makes it possible to evaluate, by means of the force model, nine accel-
erations from these positions and velocities and the position and velocity of epoch.
6.19 Integration for Elliptic Orbits of High Eccentricity 997
These accelerations must be converted from the domain of time (t) to the domain of
intermediate anomaly (s), by changing time derivatives to intermediate anomaly
derivatives. To this end, the derivatives with respect to s are computed from the
derivatives with respect to t, by means of the following relation
dð. . .Þ dð. . .Þ
¼ c1 r n
dt ds
Differentiation with respect to intermediate anomaly is indicated below by a

grave sign (`), whereas differentiation with respect to time is indicated by a prime
sign (′). The velocity vector becomes
x0 ¼ c1 r n x
and the acceleration vector becomes

0 0
x00 ¼ ðx0 Þ ¼ c1 r n x ¼ c1 r n c1 r n x
¼ c3 r 2n cx nc2 r 2n1 rxþ c2 r 2n x
An arbitrary choice of c could lead to a value of c` other than zero. Such would be
the case if c depended on orbital elements (e.g., major semi-axis or eccentricity). In
the present case, where c depends only on the gravitational parameter of the Earth
(c = l−1/2), c` is equal to zero. Therefore, x`` can be expressed as a function of x″,
as follows

x ¼ c2 r 2n x00 þ nc2 r 2n1 rx
Substituting c = l−1/2, n = 3/2, r` = c rn r′ and x` = c rn x′ in the expression written

above, we have

1 3 00 3 2 0 0
x ¼ r x þ r rx
l 2
The expression written above includes the time derivative (r′) of the radius vector
(r), which is expressible as follows
1
0 dðx xÞ2 x x0
r ¼ ¼
dt r
Hence

1 3 00 3 0 0
x ¼ r x þ r ðx x Þx
l 2
The expression written above (i.e., the second derivative of the position vector with
respect to the intermediate anomaly) is to be integrated by means of the
Gauss-Jackson and the summed Adams mid-corrector formulae to find the position
and velocity at each of the eight points around epoch. The velocity does not result
directly from the integration process, which yields x` dx/ds. Hence, the velocity
vector x′ results from the following conversion
x0 ¼ l2 r 2 x0
1 3
In order to compute the forces acting on the satellite, the time at each step must be
known. To this end, we integrate the additional equation
dt
¼ l2 r 2
1 3
t
ds
by means of the summed Adams mid-corrector formulae. Having computed the

time at each step, the forces acting on the satellite are evaluated to get refined
estimates of the accelerations (x′′). These, in turn, are converted into second
derivatives (x``) with respect to the intermediate anomaly (s). The integration is
performed again to get positions, velocities, and times at the mesh points. This
process is iterated till the accelerations between two consecutive iterations converge
to a prescribed tolerance.
This process is substantially an integration with fixed step size, with steps
equally spaced with respect to intermediate anomaly, not to time.
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Chapter 7
Dynamics of Rigid Bodies
7.1 The Motion of Rigid Bodies
A rigid body is an aggregate of particles which is not subject to deformation. In

other words, the distance between any pair of particles is constant, and the angles
formed by any triple of particles are also constant.
The position of a rigid body is defined by specifying the positions of three
non-aligned particles of the body itself. Since three co-ordinates are necessary to
define the position of a particle in space, and since the mutual distances of the
particles belonging to a rigid body are fixed by definition, then the degree of
freedom of a rigid body is 3 3 − 3 = 6. Therefore, the motion of a rigid body is
defined by a set of six quantities. Three of them are co-ordinates, which indicate the
translations of any one of its particles along the three axes X, Y, and Z of an inertial
reference system. The other three quantities are angles, which define the orientation
of the axes x, y, and z of another reference system attached to the body and moving
with respect to the inertial system XYZ. It is convenient to place the origin, O, of the
body-fixed system xyz in the centre of mass of the body. The motion of a rigid body
with respect to the inertial system XYZ is determined by specifying the position of
its centre of mass and the orientation of the body-fixed system xyz as functions of
time.

https://doi.org/10.1007/978-3-319-62220-0_7
1002 7 Dynamics of Rigid Bodies
With reference to the preceding figure, let rO be the position vector of the origin,
O, of the body-fixed system xyz with respect to the inertial system XYZ. The
orientation of the axes of the system xyz with respect to XYZ is defined by three
independent angles. These three angles plus the three components of rO represent
the six degrees of freedom of the rigid body, as has been shown above. Let us
consider an infinitesimal displacement ds of an arbitrary point P the rigid body. This
displacement may be considered as the sum of two parts. The first part is an
infinitesimal translation of the body, which translation only moves the centre of
mass of the body from an initial position to a final position, without changing the
orientation of the body-fixed axes x, y, and z. The second part is a change of
orientation of the system xyz with respect to the system XYZ. This change of
orientation is an infinitesimal rotation of the body about its centre of mass O. As a
result of these two parts of an infinitesimal displacement, the rigid body goes from
its initial position to its final position.
Following Landau and Lifchitz [1], let r OP be the position vector of the point
P mentioned above in the body-fixed system xyz. An infinitesimal displacement
ds of the point P results from an infinitesimal displacement drO of its centre of mass
and another displacement d/ r with respect to the centre of mass O, where d/ is
an infinitesimal angle of rotation, as follows
ds ¼ drO þ d/ r
By dividing all terms of the preceding equation by the infinitesimal time dt,
during which the infinitesimal displacement ds occurs, there results
v ¼ vO þ x r
where v = ds/dt, vO = drO/dt, and x = d//dt.

7.1 The Motion of Rigid Bodies 1003
In the equation written above, vO is the velocity of the centre of mass, O, of the
rigid body. This vector is also called velocity of the translational motion. The vector
x is the angular velocity of the rotational motion of the rigid body. The vectors x
and d/ are directed along the axis of rotation of the rigid body.
In other words, the velocity v of an arbitrary point P of a rigid body, with respect
to the inertial system XYZ, can be expressed as a function of the translational
velocity vO of the origin O and the rotational velocity x of the body itself.
It is to be noted that the equation v = vO + x r holds independently on
whether the origin, O, of the body-fixed system xyz be, or be not, the centre of mass
of the rigid body.
7.2 The Matrix of Inertia
Let three orthogonal axes x, y, and z be attached to a rigid body. Let the origin O of
these axes be chosen at any point. As has been shown in the preceding section, the
velocity vi of an arbitrary point Pi (i = 1, 2, …, N) of a rigid body, with respect to an
inertial system XYZ, can be expressed as follows
vi ¼ vO þ x r i
The moment of momentum of the rigid body with respect to the origin O of the
body-fixed system xyz is
X
N X
N X
N
hO ¼ ri mi ðvO þ x ri Þ ¼ ri ðx ri Þmi vO m i ri
i¼1 i¼1 i¼1
When the reference point O is in quiet with respect to XYZ, then vO = 0.

Otherwise, when O coincides with the centre of mass of the rigid body, then
X
N
mi ri ¼ 0
i¼1
Therefore, when the reference point O is either in quiet or coincident with the
centre of mass of the rigid body, then the second addend on the right side of the
expression of hO vanishes.
When the rigid body is a continuous medium instead of an aggregate of discrete
particles, then its moment of momentum is expressible as follows
Z
hO ¼ r ðx rÞdm
The vector product between parentheses is

x r ¼ xy z xz y ux þ ðxz x xx zÞuy þ xx y xy x uz
The preceding expression, multiplied by dm, yields the components along

respectively x, y, and z of the momentum vector of the particle dm.
By executing the vector product r (x r)dm, there results

r ðx rÞdm ¼ xx y2 þ z2 xy xy xz xz dm ux þ f½xx xy þ xy x2 þ z2

xz yzdmguy þ f½xx xz xy yz þ xz x2 þ y2 dmguz
The preceding equation expresses the three components along x, y, and z of the
moment of momentum vector dhO of the particle dm with respect to O.
By integrating over the whole rigid body, we obtain the components along x, y,
and z of the moment of momentum hO of the rigid body with respect to O, that is,
hO ¼ hOx ux þ hOy uy þ hOz uz
The three moments of inertia Ix, Iy, and Iz of the rigid body with respect to the
axes x, y, and z are defined as follows
R R R
Ix ¼ ðy2 þ z2 Þdm Iy ¼ ðx2 þ z2 Þdm Iz ¼ ðx2 þ y2 Þdm
The products of inertia of the rigid body are defined as follows

R R R
Ixy ¼ x y dm Ixz ¼ x z dm Iyz ¼ y z dm
As a result of this definition, the components along x, y, and z of the moment of

momentum hO of the rigid body with respect to O can be expressed as follows
hOx ¼ Ix xx Ixy xy Ixz xz

hOy ¼ Ixy xx þ Iy xy Iyz xz
hOz ¼ Ixz xx Iyz xy þ Iz xz
The moments of inertia and the products of inertia defined above may be rep-
resented by means of the following matrix of inertia =
2 3
Ix Ixy Ixz
= 4 Ixy Iy Iyz 5
Ixz Iyz Iz
7.2 The Matrix of Inertia 1005
By using the matrix of inertia defined above, the moment of momentum vector
hO of a rigid body can be expressed as follows
hO ¼ = x
7.3 Kinetic Energy of a Rigid Body
Let x, y, and z be three perpendicular axes of a body-fixed system of reference,

whose origin O is the centre of mass of a rigid body moving in space. Let P be any
other point of the body than its centre of mass, and let r be the position vector of P
in the system xyz. As has been shown in Sect. 7.1, the velocity vector of P with
respect to an inertial system XYZ can be expressed as follows
v ¼ vO þ x r
where vO is the velocity of the centre of mass of the rigid body with respect to the
The square of the velocity of P results from the following scalar product
v2 ¼ v v ¼ ðvO þ x rÞ ðvO þ x rÞ ¼ v2O þ ðx rÞ ðx rÞ þ 2vO ðx rÞ
Consequently, the total kinetic energy of the whole rigid body is

Z Z Z
1 1 1
T¼ v2 dm ¼ mv2O þ ðx rÞ ðx rÞdm þ vO x r dm
2 2 2
Z
1 2 1
¼ mvO þ ðx rÞ ðx rÞdm
2 2
R
This is because the integral r dm vanishes when the origin, O, of the system of
reference xyz coincides with the centre of mass of the rigid body.
Consequently, the total kinetic energy of a rigid body results from the sum of
two terms. Let these terms be Ttrans and Trot. The first term, Ttrans, is the kinetic
energy of the body in its motion of translation and is computed as if the whole
mass, m, of the body were concentrated in its centre of mass. The second term, Trot,
is the kinetic energy of the body in its motion of rotation about an axis (which is the
line of action of the vector x) passing through the centre of mass of the body.
Let us consider the second, Trot, of the two terms mentioned above. By executing
the scalar product (x r) (x r), there results

ðx rÞ ðx rÞ ¼ x2x y2 þ z2 þ x2y x2 þ z2 þ x2z x2 þ y2 2 xx xy xy
2 xy xz yz 2 xx xz xz
Remembering the moments and products of inertia defined in Sect. 7.2, the total
kinetic energy of a rigid body is
1 1 2
T ¼ mv2O þ xx Ix þ x2y Iy þ x2z Iz 2xx xy Ixy 2xy xz Iyz 2xx xz Ixz
2 2
where the first term on the right-hand side of the preceding equation is Ttrans, and
the second term is Trot.
7.4 Moment of Inertia of a Rigid Body About

an Arbitrary Axis
Let Inn be the moment of inertia of a rigid body about an arbitrary axis n, and let x
be the angular velocity vector of the body. As has been shown in Sect. 7.3, the
rotational part of the kinetic energy of the body can be expressed as follows
1
Trot ¼ Inn x2
2
The preceding equation, solved for Inn, yields
Trot
Inn ¼ 2
x2
1 2
Trot ¼ xx Ix þ x2y Iy þ x2z Iz 2xx xy Ixy 2xy xz Iyz 2xx xz Ixz
2
shown in Sect. 7.3, yields

x 2 x 2 x 2 x x x x
x y z x y y z
Inn ¼ Ix þ Iy þ Iz 2 Ixy 2 Iyz
x x x x x x x
xx xz
2 Ixz
x x
Let ‘nx, ‘ny, and ‘nz be the direction cosines of the vector x (or of the axis n) with
respect to the axes x, y, and z of a body-fixed system of reference, whose origin is
the centre of mass of the body. The preceding equation may also be written as
follows
Inn ¼ ‘2nx Ix þ ‘2ny Iy þ ‘2nz Iz 2‘nx ‘ny Ixy 2‘ny ‘nz Iyz 2‘nx ‘nz Ixz
7.4 Moment of Inertia of a Rigid Body About an Arbitrary Axis 1007
The preceding equation may be written in the form of a double sum, as follows
XX
Inn ¼ ‘na ‘nb Iab
a b
where a and b are indices which apply to the letters x, y, and z, so that Iaa stands for
Ixx = Ix, Iyy, Izz, whereas Iab stands for −Ixy, −Iyz, −Ixz.
Therefore, the products of inertia can be expressed as follows
XX
Ing ¼ ‘na ‘gb Iab
a b
The direction cosines ‘nx, ‘ny, and ‘nz defined above may be formed by means of
the transformation matrix between two systems of co-ordinates. Let x, y, and z be a
system of orthogonal axes, whose origin is a point O, and whose unit vectors are
respectively ux, uy, and uz. Let x*, y*, and z* be another system of orthogonal axes,
whose origin is the same point O, and whose unit vectors are respectively ux*, uy*,
and uz*. An arbitrary point P in space, whose position vector is OP r, may be
indicated as follows
r ¼ xux þ yuy þ zuz ¼ x ux þ y uy þ z uz
The transformation from the co-ordinates x*, y*, and z* to the co-ordinates x, y,
and z is expressed by the following matrix product
2 3 2 32 3
x ‘xx ‘xy ‘xz x
4 y 5 ¼ 4 ‘yx ‘yy ‘yz 54 y 5
z ‘zx ‘zy ‘zz z
where the elements of the preceding matrix are the direction cosines resulting from
the following scalar products
‘xx ¼ ux ux ; ‘xy ¼ uy ux ; ‘xz ¼ uz ux ;

‘yx ¼ uy ux ; ‘yy ¼ uy uy ; ‘yz ¼ uy uz ;
‘zx ¼ uz ux ; ‘zy ¼ uz uy ; ‘zz ¼ uz uz :
The rotational part of the kinetic energy of a rigid body may be expressed by the
scalar product of the angular velocity vector and the moment of momentum, as
follows
1
Trot ¼ ðx hO Þ
2
Since in Sect. 7.2, we have shown that hO ¼ = x, where = is the matrix of

inertia, then
1
Trot ¼ ðx = xÞ
2
Thomson [2] explains the meaning of the double scalar product x = x, as

follows. Since the matrix of inertia = can be expressed as the product of two
matrices A and B (i.e. = = AB), then the preceding equation may also be written as
follows
1
Trot ¼ ðx AÞðB xÞ
2
and therefore Trot is a scalar quantity, because it results from the product of two
scalar quantities, which are (x A) and (B x).
7.5 Principal Axes of Inertia
The principal axes of inertia of a rigid body, which are denoted here by 1, 2, and 3,
are defined as those axes along which the products of inertia vanish.
Let I1, I2, and I3 be the moments of inertia of a rigid body along its principal
axes, respectively, 1, 2, and 3. The moment of inertia of a rigid body with respect to
the instantaneous axis of rotation, n, may be expressed as a function of I1, I2, and I3
as follows
In ¼ I1 ‘2n1 þ I2 ‘2n2 þ I3 ‘2n3
where ‘n1, ‘n2, and ‘n3 are the direction cosines of the angular velocity vector x (or
of the axis n) with respect to the principal axes 1, 2, and 3. Since the axes 1, 2, and 3
are fixed to the rigid body, then the three moments of inertia I1, I2, and I3 are also
invariable. However, when the direction of the instantaneous axis of rotation, n,
changes, then the direction cosines ‘n1, ‘n2, and ‘n3 change, and therefore, the value
of In also changes. To see how the value of In changes as a function of the direction
cosines, we set q2 = 1/In, and therefore, the equation written above may also be
written as follows
I1 ð‘n1 qÞ2 þ I2 ð‘n2 qÞ2 þ I3 ð‘n3 qÞ2 ¼ 1
which is the equation of an ellipsoid, whose co-ordinates are x1 = ‘n1q, x2 = ‘n2q,

and x3 = ‘n3q, and whose principal semi-axes are a1 = (1/I1)½, a2 = (1/I2)½, and
a3 = (1/I3)½. This ellipsoid is called the ellipsoid of inertia.
7.6 Euler’s Equations 1009
7.6 Euler’s Equations
The fundamental equation governing the motion of a rigid body rotating about its
centre of mass is
d hO
MO ¼
dt
where MO is the moment of all external forces applied to a rigid body, with respect
to its centre of mass, O, and hO is the moment of momentum with respect to the
same point. This equation holds in an inertial reference system XYZ.
On the other hand, as has been shown in the preceding sections, the vector hO
has a simple expression in a system of reference whose axes are aligned with the
principal axes of inertia of the body, because, in that case, the matrix of inertia = is
a diagonal matrix. It is convenient to define two orthogonal systems of reference,
whose common origin is the centre of mass O of the body. The first of them, whose
axes are X, Y, and Z, is an inertial system. The second system, whose axes are the
principal axes of inertia 1, 2, and 3 of the body, rotates together with the body and is
therefore a non-inertial system.
Since the system 123 rotates with the body, the instantaneous angular velocity x
of this system is the same as that of the body, and therefore
d hO
h0O ¼ hO þ x hO
dt
where the apex sign (′) indicates the time derivative in the inertial system XYZ, and
the grave sign (`) indicates the time derivative in the rotating system 123. By using
this notation, the fundamental equation governing the motion of a rigid body
rotating about its centre of mass O becomes
M O ¼ hO þ x hO
Let MO1, MO2, and MO3 be the three components of the vector MO in the rotating
system 123, and let x1, x2, and x3 be the three components of the vector x in the
same system. Since the matrix of inertia = is diagonal in this system, then the three
components of the vector hO in the same system are I1x1, I2x2, and I3x3, where I1,
I2, and I3 are the moments of inertia relative to the principal axes of inertia of the
rotating body. Consequently, the vector product x hO yields
x hO ¼ ½ðI2 I3 Þx2 x3 u1 ½ðI3 I1 Þx3 x1 u2 ½ðI1 I2 Þx1 x2 u3
where u1, u2, and u3 are the unit vectors along the principal axes of inertia,
respectively, 1, 2, and 3. Since I1, I2, and I3 do not vary with time in the rotating
system, then the three scalar equations corresponding to MO = hÒ + x hO are
MO1 ¼ I1 x1 ðI2 I3 Þx2 x3

The equations written above are known as Euler’s equations. They are difficult
to solve for a rotating body of arbitrary shape.
7.7 An Axially Symmetric (I1 = I2) Rotating Body

not Subject to External Moments
Let us consider an axially symmetric body, that is, a body having two moments of
inertia equal (e.g. I1 = I2), as shown in the following figure.
The condition I1 = I2 is independent of whether the body be a cylinder, or a

parallelepiped, or a prism, or a pyramid, provided that its section perpendicular to
the axis of symmetry be a regular polygon, and the body be made of a homoge-
neous material. For example, the preceding figure shows a flat cylindrical body,
whose principal axis 3 is perpendicular to the two circular surfaces.
Let I1, I2, and I3 be the moments of inertia of the body illustrated above along its
principal axes of inertia 1, 2, and 3, such that I1 = I2. Let MO1, MO2, and MO3 be the
external moments acting about the principal axes of inertia, respectively, 1, 2, and 3
of the body itself. In case of the three external moments MO1, MO2, and MO3 being,
all of them, equal to zero, Euler’s equations become
I1 x1 þ ðI3 I1 Þx2 x3 ¼ 0

I1 x2 ðI3 I1 Þx3 x1 ¼ 0
I3 x3 ¼ 0
7.7 An Axially Symmetric (I1 = I2) Rotating Body … 1011
As a result of the third equation, x3 is constant in time in the body-fixed system.

By setting x3(I3/I1 − 1) = k, the first two equations become, respectively
x1 þ kx2 ¼ 0
x2 kx1 ¼ 0
By multiplying the first equation by x1, the second equation by x2, and adding,
there results
x1 x1 þ x2 x2 ¼ 0
which means
ðx21 þ x22 Þ ¼ 0
that is,
x21 þ x22 ¼ constant
Since x21 + x22 is constant and x3 is constant (x3¼0 ), then the magnitude
1
x ¼ x21 þ x22 þ x23 2
(or the length) of the vector x is also constant in the system attached to the body.
In addition, since MO1, MO2, and MO3 are equal to zero by hypothesis, then the
fundamental equation of rotating bodies implies dhO/dt h′O = 0, and therefore,
the moment of momentum vector hO is constant in both magnitude and direction in
an inertial system.
The two differential equations x1 þ kx2 ¼ 0 and x2 kx1 ¼ 0, with the
initial condition x1 = x12 and x2 = 0 for t = 0, have the following solution
x1 ¼ x12 cosðktÞ
x2 ¼ x12 sinðktÞ
where x12 = (x21 + x22)½ is the projection of the angular velocity vector x onto the
1,2-plane. Let a be the constant angle which x forms with the principal axis 3. The
components of x along the principal axes of inertia are
x1 ¼ x sin a cosðktÞ
x2 ¼ x sin a sinðktÞ
x3 ¼ x cos a
where k = (I3/I1 − 1)x3 = (I3/I1 − 1)x cos a. In other words, the angular velocity
vector x has a motion of precession, with respect to an inertial system, about the
principal axis 3, which is also the axis of symmetry of the rotating body shown
above, and k is the angular frequency (in radians/second) of this motion.
In this motion of precession, the angular velocity vector x sweeps out a cone,
called body cone, whose angle of semi-aperture is a, and whose vertex O is the
centre of mass of the rotating body, as shown in the following figure.
In the body-fixed system 123, the three components of the moment of

momentum vector hO along the principal axes of inertia, respectively, 1, 2, and 3
are
hO1 ¼ I1 x sin a cosðktÞ

hO2 ¼ I1 x sin a sinðktÞ
hO3 ¼ I3 x cos a
This is because the axis of symmetry 3 and the vectors x and hO lie in the same
plane. In other words, in the body-fixed system 123, the moment of momentum
vector has a motion of precession, with the same angular frequency k, about the
axis of symmetry of the body. In this motion of precession, the moment of
momentum vector sweeps out a cone, called space cone. In an inertial reference
system, the moment of momentum vector hO is fixed in space.
The space cone and the body cone are shown in the preceding figure, where the
moment of momentum vector hO has been aligned with the vertical axis Z of the
Let h be the constant angle which the moment of momentum vector, hO, forms
with the principal axis 3. The value of this angle results from
I1 x12 I1
tan h ¼ ¼ tan a
I3 x3 I3
This is because x12 = x sin a, and x3 = x cos a. Consequently, if I3 > I1, then
a is greater than h, as shown on the left-hand side of the preceding figure, and
therefore, hO lies between axis 3 and x. Such is the case with an oblate (i.e. flat)
body, for example, with a disc.
By contrast, if I3 < I1, then a is less than h, as shown on the right-hand side of
the preceding figure, and therefore, x lies between axis 3 and hO. Such is the case
with a prolate (i.e. elongate) body, for example, with a slender rod.
In both cases, the motion of precession can be described as the rolling, without
slipping, of an imaginary moving cone (attached to the body) over the surface of an
immovable space cone.
The equations written above lead to the following conclusions:
• the two vectors x and hO and the axis of symmetry 3 of the body lie in the same
plane, because these vectors have a motion of precession, in the body-fixed
system, with the same constant angular frequency k, about this axis (in other
words, the plane indicated above rotates at the angular velocity k);
• in the plane indicated above, the angle h (between 3 and hO) and the angle a
(between 3 and x) are constant in time;
• the two vectors hO and x are constant in magnitude; and
• the vector hO is constant in both magnitude and direction in an inertial system of
reference.
7.8 An Axially Symmetric (I1 = I2) Rotating Body

not Subject to External Moments (in Terms of Euler’s
Angles)
The present section considers the same problem as that of Sect. 7.7, in terms of the
classical Eulerian angles w, h, and /. We use the rotation sequence (3, 1, 3), also
known as the x-convention, which is described below and can also be found in [3].
These angles are three independent quantities which define the orientation a
body-fixed system with respect to an inertial system.
Following the notation used in the preceding section, we consider the body-fixed
principal axes of inertia 123 and the axes XYZ of an inertial system. The two
systems of reference indicated above have their common origin in the centre of
mass, O, of the rotating body. The straight line ON, called line of nodes, results
from the intersection of the plane XY with the plane 12. The axes 123 are supposed
to have the same directions as those of the axes XYZ at the initial time.
According to the rotation sequence (3, 1, 3), the inertial system XYZ is rotated
three times: first about the Z-axis by an angle w; then about the new X-axis (which
is the line ON) by an angle h; and finally about the newest Z-axis (which is the
3-axis) by an angle /. When the three rotations are executed in the correct
sequence, the rotated system 123 has the orientation shown in the following figure,
which is due, mutatis mutandis, to the courtesy of NASA [4].
Let R3(w), R1(h), and R3(/) be the matrices which define the three rotations
(respectively, w, h, and /) indicated above. The components of these matrices are
2 3 2 3
cos w sin w 0 1 0 0
R3 4 sin w cos w 0 5 R1 4 0 cos h sin h 5
2 0 0 13 0 sin h cos h
cos / sin / 0
R3 4 sin / cos / 0 5
0 0 1
The rotation matrix R = R3(w) R1(h) R3(/), which defines the transformation
2 3 2 3
1 X
4 2 5 ¼ R4 Y 5
3 Z
from the co-ordinates XYZ to the co-ordinates 123, has the following components
cos w cos / sin w cos h sin / cos w sin / sin w cos h cos / sin w sin h
sin w cos / þ cos w cos h sin / sin w sin / þ cos w cos h cos / cos w sin h
sin h sin / sin h cos / cos h
Since the matrix R indicated above is orthogonal (R−1 = RT), then the inverse
transformation (from 123 to XYZ) is performed through the transpose matrix RT.
The transformation of the rates w`, h`, and /` of the Eulerian angles w, h, and /
into the components x1, x2, and x3 of the angular velocity vector x is performed
as follows. Since h′O = 0, then it is convenient to orient the constant vector hO
along the Z-axis, and therefore, w` is the angular velocity of the line of nodes and
the spin axis about the Z-axis. By inspection of the preceding figure, we have
x1 ¼ wsin h sin / þ h cos /

x2 ¼ wsin h cos / h sin /
x3 ¼ wcos h þ /
Since Sect. 7.7 has shown that h is a constant angle, then h` = 0. In this case, the
preceding expressions become
x1 ¼ wsin h sin /
x2 ¼ wsin h cos /
x3 ¼ wcos h þ /
where x3 is constant in time, as has also been shown in Sect. 7.7.

The angular accelerations result from the time derivatives (where h` = 0) of the
preceding expressions. This yields
x1 ¼ w/ sin h cos /

x2 ¼ w/ sin h sin /
x3 ¼ 0
Substituting these expressions into the first of the following equations
I1 x1 þ ðI3 I1 Þx2 x3 ¼ 0

I1 x2 ðI3 I1 Þx3 x1 ¼ 0
I3 x3 ¼ 0
which have been derived in Sect. 7.7, there results
I1 w/ sin h cos / þ ðI3 I1 Þðw2 sin h cos h cos / þ w/ sin h cos /Þ ¼ 0
Simplifying and solving the preceding equation with respect to w`, there results
I3 /
w ¼
ðI3 I1 Þ cos h
The equation written above means that the roll axis rotates about the fixed Z-axis
(or about the vector hO) at a speed w` which is proportional to the angular velocity
of spin /`. The minus sign in front of I3/` means that, if I3 > I1 (as is the case with
an oblate body), then the speed is negative, because w` points in the positive
direction of the Z-axis, as shown in the second figure of Sect. 7.7. In this case, the
motion of spin is opposed to the motion of precession, and therefore, the motion of
precession is regressive (or retrograde). By contrast, if I3 < I1 (as is the case with a
prolate body), then /` and w` point in the same direction, and therefore, the motion
of precession is progressive (or prograde).
The rotating body shown in the figures of Sect. 7.7 is a flat disc, which spins
about an axis perpendicular to its circular surfaces, and has therefore a retrograde
motion of precession. If the rotating body had been a slender rod spinning about its
longitudinal axis, then its motion of precession would have been prograde.
In addition, if I3 > I1, then the ratio I3/[(I3 − I1) cos h] is greater than unity. This
means that w` is greater than /`. As h approaches 90°, w` becomes larger and larger
in comparison with /`. Since h` = 0, then the angular velocity vector x results from
the vector sum of w` and /`, as shown in the following figure.
The cones of precession are those described in Sect. 7.7.
7.9 Unsymmetrical Body Not Subject to External

Moments (Geometric Solution)
In the general case of a rotating body without a rotational symmetry and with zero
external moments applied to it, Euler’s equation (see Sect. 7.6) are
I1 x1 ðI2 I3 Þx2 x3 ¼ 0

I2 x2 ðI3 I1 Þx3 x1 ¼ 0
I3 x3 ðI1 I2 Þx1 x2 ¼ 0
where 1, 2, and 3 are the axes of the body-fixed system of reference and also the
principal axes of inertia of the body. The present section shows a geometric solution
of these equations, which is due to Poinsot [5]. The analytic solution will be shown
in Sect. 7.10.
Since no external moments are applied to the body and no work is done on it, the
kinetic energy, Trot, and the moment of momentum, hO, of the body about its centre
of mass, O, are two constant quantities. The angular velocity, x, and the moment of
momentum, hO, of the body have the following components
7.9 Unsymmetrical Body Not Subject to External Moments … 1017
x ¼ x1 u1 þ x2 u2 þ x3 u3
hO ¼ I1 x1 u1 þ I2 x2 u2 þ I3 x3 u3
along u1, u2, and u3, which are the unit vectors along the principal axes of inertia,
respectively, 1, 2, and 3 of the body.
The scalar product of the two vectors x and hO yields
x hO ¼ I1 x21 þ I2 x22 þ I3 x23 ¼ 2Trot ¼ constant
The magnitude of hO is
1 1
hO ¼ ðhO hO Þ2 ¼ I12 x21 þ I22 x22 þ I32 x23 2 ¼ constant
Let
hO
uh ¼
hO
be the unit vector directed along hO.
With reference to the preceding figure, the component of x along hO results from
2Trot
ON ¼ x uh ¼ ¼ constant
hO
The tip of the vector x lies in a plane passing through N and perpendicular to the
direction, ON, of the vector hO. Therefore, ON is called the invariable line, and the
plane perpendicular to ON is called the invariable plane. The tip of the vector x
moves in the invariable plane.
Let x1, x2, and x3 be the three components of x along the three principal axes
of inertia 1, 2, and 3 of the body, which means that x is expressed as follows
x ¼ x1 u1 þ x2 u2 þ x3 u3
where u1, u2, and u3 are the unit vectors along the principal axes of inertia of the
body. We set
1 1 1
2Trot 2 2Trot 2 2Trot 2
a1 ¼ a2 ¼ a3 ¼
I1 I2 I3
Therefore, the preceding equation
x21 x22 x23

þ þ ¼1
a2
1 a2
2 a2
3
This is the equation of an ellipsoid, called the Poinsot ellipsoid, which is shown
in the following figure.
The Poinsot ellipsoid rolls on the invariable plane so that the distance ON from
the centre of the ellipsoid to the invariable plane is
2Trot
ON ¼
hO
To prove this, let us consider again the equation
x hO ¼ 2Trot ¼ constant
Since hO and Trot are, both of them, constant quantities, then the preceding
expression can be differentiated with respect to time, as follows
dð2Trot Þ ¼ dx hO ¼ 0
The vanishing scalar product shows that the angle between dx and hO is equal to
90°, and therefore, dx and hO are mutually perpendicular.
Since the tip of x moves on the invariable plane, then any variation dx of x is
perpendicular to hO.
In addition, since the tip of x lies on the surface of the Poinsot ellipsoid, then
this ellipsoid is tangent to the invariable plane.
Consequently, the motion of the body, in an inertial system of reference, is
described by the rolling of the Poinsot ellipsoid on the invariable plane.
By contrast, in a system of reference attached to the rotating body, the invariable
plane moves with respect to the body. The equation written above
1 1
x21 x22 x23

2 þ 2 þ 2 ¼ 1
hO hO hO
I1 I2 I3
The surface described by the preceding equation is an ellipsoid, whose principal

semi-axes are a^1 = hO/I1, a^2 = hO/I2, and a^3 = hO/I3. This surface is called the
moment-of-momentum ellipsoid.
The curve traced by the tip of the vector x is defined by the intersection between
the Poinsot ellipsoid and the moment-of-momentum ellipsoid.
The instantaneous line of action of x passes through the point of contact
between the Poinsot ellipsoid and the invariable plane. Therefore, this line generates
simultaneously two cones: one in the fixed space and the other in the space attached
to the body (or in the Poinsot ellipsoid). The cone traced out in the inertial space is
called the herpolhode cone, and the cone traced out in the body-fixed space is called
the polhode cone [2]. These cones are shown in the following figure. The Poinsot
ellipsoid rolls, without slipping, on the invariable plane, so that the centre of the
Poinsot ellipsoid is at a constant height above the invariable plane. The curve traced
out by the point of contact on the Poinsot ellipsoid is the polhode, and the similar
curve on the invariable plane is the herpolhode.
The Poinsot ellipsoid and the moment-of-momentum ellipsoid have the vector x
in common. Consequently, the equation of the polhode cone can be obtained by
subtracting the equation of the Poinsot ellipsoid from the equation of the
moment-of-momentum ellipsoid and multiplying by h2O. This yields

h2 h2 h2
I1 I1 O x21 þ I2 I2 O x22 þ I3 I3 O x23 ¼ 0
2Trot 2Trot 2Trot
The preceding equation shows that, in order for the polhode cone to be real, the
quantity h2OTrot must lie between the least and the greatest of the values of the
moments of inertia I1, I2, and I3. For example, let us suppose that I3 < I2 < I1. The
condition indicated above can be expressed as follows
I3 h2O Trot I1
In order to determine the polhode curves, let us consider the intersection between
the Poinsot ellipsoid and the polhode cone. As has been shown above, the equations
of these two surfaces are, respectively

I1 I2 I3
x1 þ
2
x2 þ
2
x23 ¼ 1
2Trot 2Trot 2Trot

h2O h2O h2O
I1 I1 x1 þ I2 I2
2
x2 þ I3 I3
2
x23 ¼ 0
2Trot 2Trot 2Trot
By eliminating in sequence one of the co-ordinates x1, x2, and x3 between the
equations written above, there results

h2O
I1 ðI1 I3 Þx21 þ I2 ðI2 I3 Þx22 ¼ 2Trot I3
2Trot
2
hO
I1 ðI1 I2 Þx1 I3 ðI2 I3 Þx3 ¼ 2Trot
2 2
I2
2Trot

h2
I2 ðI1 I2 Þx22 þ I3 ðI1 I3 Þx23 ¼ 2Trot I1 O
2Trot
If ½h2O/Trot = I3, then the first equation can be satisfied for x1 = x2 = 0. In this
case, the polhode curve degenerates to a point on the x3-axis.
If ½h2O/Trot = I2, then the second of the three equations yields
1
x3 I1 ðI1 I2 Þ 2
¼
x1 I3 ðI2 I3 Þ
which indicates two planes passing through the x2-axis. The projections onto the
x1x2-plane and onto the x2x3-plane from the other two equations are ellipses.
If ½h2O/Trot = I1, then the third of the three preceding equations can be satisfied
only if x2 = x3 = 0, and the polhode curve degenerates to a point on the x1-axis.
If ½h2O/Trot lies between I2 and I3, then the polhodes lie between the planes of the
equation written above x3/x1 = {[I1(I1 − I2)]/[I3(I2 − I3]}½ and the x3-axis. Their
projections onto the x1x2-plane are ellipses.
If ½h2O/Trot lies between I1 and I2, then the polhodes lie in the central part of the
Poinsot ellipsoid between the planes defined by the equation written above
x3/x1 = {[I1(I1 − I2)]/[I3(I2 − I3]}½. Their projections onto the x2x3-plane are
ellipses, and their projections onto the x1x3-plane are hyperbolas. A drawing of the
polhodes can be found in [2].
The two equations derived above

1 1
may be rewritten in terms of the components hO1, hO2, and hO3 of the
moment-of-momentum vector hO along the three principal axes of inertia, as follows
h2O1 h2O2 h2O3

x hO ¼ þ þ ¼ 2Trot
I1 I2 I3
h2O1 þ h2O2 þ h2O3 ¼ h2O
The two equations written above represent, respectively, the Poinsot ellipsoid
and a sphere, which is called the sphere of the moment of momentum. The solution
of these equations is the intersection between the ellipsoid and the sphere.
The equations derived above

h2O h2O h2O
I1 I1 x1 þ I2 I2
2
x2 þ I3 I3
2
x23 ¼ 0
2Trot 2Trot 2Trot
2
hO
I1 ðI1 I3 Þx1 þ I2 ðI2 I3 Þx2 ¼ 2Trot
2 2
I3
2Trot
2
hO
I1 ðI1 I2 Þx1 I3 ðI2 I3 Þx3 ¼ 2Trot
2 2
I2
2Trot

h2O
I2 ðI1 I2 Þx2 þ I3 ðI1 I3 Þx3 ¼ 2Trot I1
2 2
2Trot
may be rewritten in terms of the components hO1, hO2, and hO3 of the
moment-of-momentum vector hO as follows

h2O h2O h2O
1 h þ 1
2
h þ 1
2
h2 ¼ 0
2I1 Trot O1 2I2 Trot O2 2I3 Trot O3
2
I1 I3 2 I2 I3 2 hO
hO1 þ hO2 ¼ 2Trot I3
I1 I2 2Trot
2
I1 I2 2 I2 I3 2 hO
hO1 hO3 ¼ 2Trot I2
I1 I3 2Trot

I1 I2 2 I1 I3 2 h2
hO2 þ hO3 ¼ 2Trot I1 O
I2 I3 2Trot
The solution of these equations is the intersection between the Poinsot ellipsoid and
the sphere of the moment of momentum. In the particular case, described in Sect. 7.7, of
an axially symmetric rotating body (I1 = I2) with zero external moments, the Poinsot
ellipsoid is an ellipsoid of revolution, and consequently, the polhode is a circumference
about the axis of symmetry of the body. Likewise, the herpolhode on the invariable
plane is a circumference. An observer in the system of reference attached to the body
sees the angular velocity vector x move on the surface of the space cone, whose
intersection with the invariable plane is the herpolhode. Therefore, the free motion of an
axially symmetric rigid body may be described in geometric terms as the rolling of the
body cone on the space cone, as has been shown in Sect. 7.7.
7.10 Unsymmetrical Body Not Subject to External

Moments (Analytic Solution)
In case of a rotating body whose moments of inertia along the principal axes are not
equal, the analytic solution of the Euler equation requires the use of elliptic func-
tions. Starting from the two following equations derived in Sect. 7.9

1 1
and eliminating in sequence x1, x2, and x3, there results
I2 ðI2 I1 Þx22 þ I3 ðI3 I1 Þx23 ¼ h2O 2I1 Trot

By solving the third and the first of the three equations written above for,
respectively, x21 and x23, there results

h2O 2I3 Trot I2 ðI2 I3 Þ 2
x21 ¼ 1 2 x2
I1 ðI1 I3 Þ hO 2I3 Trot
" 2 #
2I1 Trot hO2
I1 I2 hO 2I3 Trot I2 ðI2 I3 Þ
x3 ¼
2
1 x 2
I3 ðI1 I3 Þ I2 I3 2I1 Trot h2O h2O 2I3 Trot 2
The two equations written above can be simplified by setting

" #12
I2 ðI2 I3 Þ
y¼ 2 x2
hO 2I3 Trot
2 1
I1 I2 hO 2I3 Trot 2
k¼
I2 I3 2I1 Trot h2O
where 0 k 1. By so doing, there results

1
h2O 2I3 Trot 2 1
x1 ¼ 1 y2 2
I1 ðI1 I3 Þ
1
2I1 Trot h2O 2 1
x3 ¼ 1 k 2 y2 2
I3 ðI1 I3 Þ
The preceding expressions, substituted into
I2 x2 ðI3 I1 Þx3 x1 ¼ 0
(which is the second of the three Euler equations of Sect. 7.6), yield
1
h2O 2I3 Trot 2
x2 ¼ y
I2 ðI2 I3 Þ
" #12
I3 I1 h2O 2I3 Trot 2I1 Trot h2O 1
¼ 2
1 y2 1 k 2 y2 2
I2 I1 I3 ðI3 I1 Þ
The preceding expression of x`2 is an elliptic integral of the first kind, that is,
1 Zy
ðI2 I3 Þ 2I1 Trot h2O 2 dy
u ¼ Nt ¼ t¼ 1
I1 I2 I3 ½ ð 1 y2 Þ ð 1 k 2 y2 Þ 2
0
where time, t, is measured from the moment in which x2 = 0.

By setting y = sin /, there results
Z/
dh
u¼ 1 ¼ F ð/; kÞ
0
1 k 2 sin2 h 2
Hence, u is a function of the amplitude, /, and the elliptic modulus, k.

Inversely, y is a function of u = Nt and k and can be found as a tabulated
function for 0 k 1, where N is defined by the preceding equation
1 Zy
ðI2 I3 Þ 2I1 Trot h2O 2 dy
u ¼ Nt ¼ t¼ 1
I1 I2 I3 ½ ð 1 y2 Þ ð 1 k 2 y2 Þ 2
0
As is the case with the direct circular function y = sin(x) and the corresponding
inverse function x = arcsin(y), so here the dependency of y on u and k is expressed
as follows
y ¼ snðu; k Þ
Therefore, the solution of the Euler equations for the component x2 of the
angular velocity vector x along the principal axis 2 can be expressed as follows
12
h2 2I3 Trot
x2 ¼ O snðNt; kÞ
I2 ðI2 I3 Þ
Likewise the solution of the same equations for the other two components,
respectively, x1 and x3, can be expressed as follows
2 1
hO 2I3 Trot 2
x1 ¼ cnðNt; kÞ
I1 ðI1 I3 Þ
1
2I1 Trot h2O 2
x3 ¼ dnðNt; kÞ
I3 ðI1 I3 Þ
The functions cn and dn are related to the function sn as follows
cn2 x ¼ 1 sn2 x
dn2 x ¼ 1 k 2 sn2 x
The solutions given above correspond to the case in which h2O < 2I2Trot, which
condition is required for 0 k 1.
When the value of the modulus k approaches zero, that is, when I1 I2, then the
elliptic functions approach the corresponding circular functions, as follows
snðNt; 0Þ ¼ sinðNtÞ
cnðNt; 0Þ ¼ cosðNtÞ
dnðNt; 0Þ ¼ 1:0
Consequently, when I1 = I2, then the three components of the angular velocity
vector x along the three principal axes of inertia are
x1 ¼ x sin a cosðktÞ
x2 ¼ x sin a sinðktÞ
x3 ¼ x cos a
where k = (I3/I1−1)x3 = (I3/I1−1)x cos a, and a is the constant angle which the
angular velocity vector x forms with the principal axis 3, in accordance with the
results found in Sect. 7.7.
7.11 Elementary Concepts on Elliptic Integrals
The adjective elliptic used with integrals arises from the problem of measuring the
length of an arbitrary arc of ellipse. Elliptic integrals are to an ellipse what inverse
trigonometric functions are to a circumference, as the sequel will show.
Let us consider the well-known equation of an ellipse, whose major and minor
semi-axes are, respectively, a and b, in an orthogonal system of reference xy, whose
origin, O, is the centre of the ellipse:
x2 y2
þ ¼1
a2 b2
With reference to the equation written above, let us consider, without loss of
generality, an elliptic arc whose origin is on the y-axis (x = 0) and whose length, s,
is less than or equal to an elliptic quadrant:
Zx Zx 1
12 a2 ð1 b2 =a2 Þx2 2
s¼ dx þ dy
2 2
¼ dx
a2 x 2
0 0
After setting for convenience k2 = 1 − b2/a2, where 0 k 1 is the eccen-

tricity of the given ellipse, and operating a change of variable from x to t = x/a, the
preceding expression may be written as follows
Za 12
x
1 k2 t2
s¼a dt
1 t2
0
By setting t = sin /, the preceding expression may also be written as follows
Z/
1
s¼a 1 k2 sin2 h 2 dh
0
The integral on the right-hand side of the two preceding expressions is also
denoted by E(x, k). It is a function of the two variables x and k (or / and k) and is
called incomplete elliptic integral of the second kind. In case of the limits of
integrations being 0 and //2, the same integral is called complete integral of the
second kind. There are two further kinds of elliptic integral. They are:
Z/
dh
u ¼ F ð/; kÞ ¼ 1
0
1 k 2 sin2 h 2
which is the incomplete elliptic integral of the first kind, introduced to the reader in
Sect. 7.10, and
Z/
dh
Pðn; /; kÞ ¼ 1
0
1 n sin h 1 k 2 sin2 h 2
2
which is the incomplete elliptic integral of the third kind, and n is a constant called
elliptic characteristic.
7.11 Elementary Concepts on Elliptic Integrals 1027
The analogy existing between inverse circular functions and elliptic integrals
also holds with direct circular functions and elliptic functions. To show this, let us
consider the following function
1
f ðt Þ ¼ 1
ð1 t 2 Þ2
As is well known, the integral between 0 and x of the function written above is
Zx
f ðtÞdt ¼ arcsinð xÞ
0
where the function arcsin(x) is defined in the domain −1 x 1, and its range is
restricted to the principal branch −p/2 arcsin(x) p/2. The first elliptic
function, sn, is the inverse function of u = F(/, k) with respect to x, that is,
x ¼ snðu; kÞ ¼ sin /
As is the case with the trigonometric cosine, so there exists a corresponding

elliptic function cn(u, k) such that
1 1
cnðu; kÞ ¼ 1 sn2 ðu; kÞ 2 ¼ 1 x2 2 ¼ cos /
The third elliptic function dn(u, k) is defined as follows

1 1
dnðu; kÞ ¼ 1 k 2 x2 2 ¼ 1 k 2 sin2 / 2 ¼ Dð/Þ
When the value of the modulus k tends to zero, then sn(u, k) and cn(u, k) tend to,
respectively, sin u and cos u, and dn(u, k) tends to unity, as follows
snðu; 0Þ ¼ sin u cnðu; 0Þ ¼ cos u dnðu; 0Þ ¼ 1
When the value of k tends to unity, then sn(u, k) tends to tanh u, and cn(u, k) and
dn(u, k) tend both of them to 1/cosh u, as follows
snðu; 1Þ ¼ tanh u cnðu; 1Þ ¼ cosh

1
u dnðu; 1Þ ¼ cosh u
1
The upper limit of integration, /, in the elliptic integral of the first kind
Z/
dh
u ¼ F ð/; kÞ ¼ 12
0
1 k2 sin2 h
is called the Jacobi amplitude. It is a function of u, and is denoted by am u. In terms

of sn(u, k), the Jacobi amplitude, /, is defined by sin / = sn(u, k).
The figure shown below, due to Greenhill [6], illustrates the following elliptic
functions: (i) y = am x, (ii) y = cn x, (iii) y = sn x, and (iv) y = dn x.
The value of k is equal to sin 45° 0.7071. The quantity K is called the quarter
period of the elliptic functions, to the modulus k. The quarter period is defined by
p
Z2
d/
K¼ 12
0
1 k2 sin2 /
As shown in the preceding figure, the curves y = sn x and y = cn x are essen-

tially distinct curves, which cannot be superposed, like y = sin x and y = cos x.
The curve (i), which is the graph of am x, is a regular undulation running along
the straight line y = ½px/K.
The elliptic integrals arise not only in physical problems, such as those con-
cerning the motion of a rigid body or a pendulum, but also in mathematical
problems, because they are the solutions of some differential equations, as will be
shown below. Let us consider the following differential equation
d2 x
x00 ¼ x
dt2
As is well known, the solution of this equation is a combination of trigonometric

sines and cosines. By multiplying both terms of the preceding equation by x′, there
results
0
x00 x0 þ xx0 ¼ ½ðx0 Þ2 0 þ x2 ¼ 0
which, integrated once with respect to t, yields
ðx0 Þ2 þ x2 ¼ constant ¼ 1
having chosen unity for the arbitrary constant of integration, or
ðx0 Þ2 ¼ 1 x2
The preceding equation, integrated again with respect to t, yields
arcsinðx x0 Þ ¼ t t0
where the function on the left-hand side is an inverse circular function.

If, instead of (x′)2 = 1 − x2, we consider the following differential equation
ðx0 Þ2 ¼ ð1 x2 Þð1 k2 x2 Þ
where k is a constant, then integration with respect to t yields

Zx
dx
t t0 ¼ 1
x0
½ð1 x2 Þð1 k2 x2 Þ2
If we impose the condition x = 0 for t = 0, then the preceding expression

becomes
Zx
dx
t¼ 1
½ð1 x2 Þð1 k 2 x2 Þ2
0
which is the incomplete elliptic integral of the first kind, as has been shown above.
Let us consider now the following system of differential equations
x01 ¼ x2 x3
x02 ¼ x1 x3
x03 ¼ k 2 x1 x2
where one of the initial conditions prescribes x2 = 0 at t = 0.

This system of differential equations has the following solution
x1 ¼ cnðt; kÞ
x2 ¼ snðt; k Þ
x3 ¼ dnðt; k Þ
To show this, we take
x23 ¼ 1 k2 x22
x21 ¼ 1 x22
The first of these two equations, differentiated with respect to t, yields
x03 x3 ¼ k 2 x02 x2
which satisfies the given system of differential equations. This is because the last
equation written above can be obtained as a result of the following operations:
(a) the second equation (x′2 = x1x3) of the given system is multiplied by k2x2:
k2 x02 x2 ¼ k2 x1 x2 x3
(b) the third equation (x′3 = −k2x1x2) of the given system is multiplied by x3:
x03 x3 ¼ k 2 x1 x2 x3
(c) the terms of the last two equations written above are summed one to another.
Now, the second (x21 = 1 − x22) of the two equations
x23 ¼ 1 k2 x22
x21 ¼ 1 x22
written above is differentiated with respect to t. This yields
x01 x1 ¼ x02 x2
which also satisfies the given system of differential equations.

This is because the last equation written above can be obtained as a result of the
following operations:
(a) the first equation (x′1 = −x2x3) of the given system is multiplied by x1:
x01 x1 ¼ x1 x2 x3
(b) the second equation (x′2 = x1x3) of the given system is multiplied by x2:
x02 x2 ¼ x1 x2 x3
(c) the terms of the last two equations written above are summed one to another.
Finally, by substituting the preceding expressions
x23 ¼ 1 k2 x22
x21 ¼ 1 x22
into the square of the second equation (x′2 = x1x3) of the given system, there results

ðx02 Þ2 ¼ ðx1 Þ2 x23 ¼ 1 x22 1 k2 x22
which is just the differential equation
ðx0 Þ2 ¼ ð1 x2 Þð1 k2 x2 Þ
considered above. By taking account of the initial condition x2 = 0 for t = 0, the

three following functions
x1 ¼ cnðt; kÞ x2 ¼ snðt; kÞ x3 ¼ dnðt; kÞ
are the solution of the given system of differential equations.

Of course, the other two initial conditions must be satisfied. By taking account of
x23 ¼ 1 k2 x22
x21 ¼ 1 x22
these two conditions are x1 = 1 and x3 = 1 at t = 0. These conditions prescribe the

values of the functions cn(t, k) and dn(t, k) at t = 0. The other condition (x2 = 0 at
t = 0) considered above prescribes the value of the function sn(t, k) at t = 0.
7.12 Stability of the Rotation of a Rigid Body About Its

Principal Axes
By stability of a motion, we mean the sensitivity, or rather the insensitivity, of that

motion to small perturbations. Let us write Euler’s equations in case of a freely
rotating rigid body. In these conditions, Euler’s equations given in Sect. 7.6 are the
following
I1 x1 ¼ ðI2 I3 Þx2 x3

I2 x2 ¼ ðI3 I1 Þx3 x1
I3 x3 ¼ ðI1 I2 Þx1 x2
where 1, 2, and 3 are the axes of the body-fixed system of reference and also the
principal axes of inertia of the body. The equations written above lead to the
following results:
if x2 ¼ x3 ¼ 0; then x1 ¼ constant
In other words, permanent rotations are possible about each of the principal axes
of inertia of the body. As will be shown below, these permanent rotations are stable
about the axes of minimum and maximum moment of inertia and unstable about the
axis of intermediate moment of inertia of the body.
To prove this, we assume a constant rotation about one of the principal axes of
inertia (e.g. about the principal axis 1) and the presence of a small perturbation. The
initial conditions (t = 0) in the absence of the perturbation are:
x1 ¼ constant ¼ x0 ; and x2 ¼ x3 ¼ 0
The perturbing condition is
x1 ¼ x0 þ e; such that x2 and x3 are small with respect to x1
The Euler equations relating to this case can be rewritten as follows

x1 ¼ 0
I2 x2 ¼ ðI3 I1 Þx3 x1
I3 x3 ¼ ðI1 I2 Þx1 x2
By differentiating the second equation with respect to time and substituting the
value of x`1 coming from the first equation and the value of x`3 coming from the
third equation, there results
ðI1 I2 ÞðI1 I3 Þ
x2 þ ¼ x20 x2
I2 I3
Likewise, by differentiating the third equation with respect to time and substi-
tuting the value of x`1 coming the first equation and the value of x`2 coming from
the second equation, there results
ðI1 I2 ÞðI1 I3 Þ
x3 þ ¼ x20 x3
I2 I3
Therefore, in order for the constant rotation of the rigid body about the principal
axis 1 to be stable, the value of the product (I1 − I2)(I1 − I3) must be positive.
7.12 Stability of the Rotation of a Rigid Body About Its Principal Axes 1033
This happens either when I1 > I2 and I1 > I3, or when I1 < I2 and I1 < I3.
In the first case (I1 > I2 > I3), the stable rotation occurs about the axis of the
maximum moment of inertia. In the second case (I1 < I2 < I3), the stable rotation
occurs about the axis of the minimum moment of inertia. When I2 < I1 < I3, then
the product (I1 − I2)(I1 − I3) has negative values, and therefore, small values of x2
and x3 tend to increase.
The first (x`1 = 0) of the three Euler equations written above implies the con-
stancy in time of the small perturbation e. Therefore, when the principal axis 1 of
the body is either the axis of maximum or the axis of minimum moment of inertia,
then the rotation of the body about its principal axis 1 is stable. By contrast, when
this condition is not satisfied (I2 < I1 < I3), then small values of x2 and x3 tend to
increase, and the motion is unstable.
7.13 General Motion of a Rigid Body
In the cases considered in the preceding sections, the two following equations
f ¼ mðvO þ x vO Þ
M O ¼ hO þ x hO
which describe, respectively, the translational motion of the centre of mass O and
the rotational motion of a rigid body about its centre of mass or about a fixed point
can be solved independently of each other.
There are more general types of motion of a rigid body in which the force f may
contribute to the moment MO, and therefore, the two equations written above are
coupled. By considering the 3 + 3 = 6 scalar components of the force and moment
indicated above, we have six interdependent equations which are to be solved
simultaneously. The problems arising in these cases are more difficult to solve than
are those which have been considered previously. One of these problems, proposed
by Thomson [2], is considered below:
A thin circular disc rolling on a rough horizontal plane.
Let us consider the disc shown in the following figure. Let XYZ be a system of
reference whose axes X, Y, and Z have fixed directions, and whose origin O
coincides with the centre of mass of the disc. Let xyz be another system of reference
having its origin in O and moving with the disc.
The z-axis is perpendicular to the plane of the disc. The axes x and y are in the
plane of the disc, so that x is perpendicular to the plane of the sheet.
Let us consider a further axis, η, such that x, η, and Z form a third set of axes. This
set is rotated from the X-, Y-, and Z-axes about the Z-axis. As shown in the preceding
figure, the η-axis is the projection of the y-axis onto the horizontal direction. The
orientations of the axes z and x are defined by a rotation h about the x-axis for z and a
rotation w about the Z-axis for x. The floor exerts a constraint force on the disc. This
force has a perpendicular component, fZ, directed upward along the Z-axis and a
friction component on the floor. This friction component, in turn, has a component, fx,
directed along x, and another component, fη, directed along η.
Let I1, I1, and I3 (with I1 < I3) be the moments of inertia of the disc about,
respectively, x, y, and z. Let R be the radius of the disc. The three scalar equations
expressing the moments MOx, MOy, and MOz about, respectively, x, y, and z are
MOx ¼ I1 h I3 xz w sin h I1 w2 sin h cos h ¼ fZ R cos h þ fg R sin h

MOy ¼ I1 ðw sin h þ wh cos hÞ þ I1 hw cos h þ I3 xz h ¼ 0
MOz ¼ I3 xz ¼ fx R
where xz = /` − w` cos h in the negative direction of the z-axis.

The second equation, which expresses MOy, may be rewritten as follows
I1 wsin h ¼ hðI3 xz þ 2I1 w cos hÞ

According to the preceding equation, in case of h` being negative (i.e. in case of
the disc falling), then w`` increases, or the spin about the vertical Z-axis increases,
and the disc rolls into a smaller circle. Since the disc rolls without slipping on the
7.13 General Motion of a Rigid Body 1035
floor due to the friction, then the velocity vector, vO, of its centre of mass has the
following components along the axes defined above
vOx ¼ Rxz
vOz ¼ Rh cos h
vOZ ¼ Rh
vOg ¼ Rh sin h
These components are shown in the following figure.
The three scalar equations expressing the forces along, respectively, x, η, and
Z are
fx ¼ mðRxz þ Rhw sin hÞ

fg ¼ mðRh sin h þ Rh2 cos h Rxz wÞ
fZ mg ¼ mðRh cos h Rh2 sin hÞ
The three equations expressing the components of the moment vector MO and
the three equations expressing the components of the force vector f acting on the
disc must be solved simultaneously. The solution of these equations is simple only
in some cases, which are indicated below.
(a) Disc rolling in a nearly vertical plane (h p/2) and xz large

In this case, the angular velocities h` and w` are negligible in comparison with xz.
In addition, since h p/2, then sin h 1. We define a further variable, that is, the
angle a such that a = p/2 − h, whose value is nearly equal to zero.
By differentiating once and twice the equation a = p/2 − h with respect to time,
there results a` = −h` and a`` = −h``.
Therefore, −a` can be written in place of h, and −a`` can be written in place of
h`` into the six scalar equations which express the components of the moment
vector MO and the components of the force vector f acting on the disc.
After a change of variable (a in place of h), these equations become, respectively
I1 a I3 xz w ¼ fZ Ra þ fg R

I1 w ¼ I3 xz a
fx R ¼ I3 xz
fx ¼ mRxz
fg ¼ mRða þ xz wÞ
fZ mg ¼ 0
The equations fx R = I3 x`z and −fx = m R x`z imply
ðI3 þ mR2 Þxz ¼ 0
hence
xz ¼ constant ¼ n
The initial conditions are w`0 = 0 and a0 = 0 at t = 0. In other words, the disc
starts in a straight path and in a vertical plane.
By integrating the second of these differential equations
I1 w ¼ I3 xz a
with the initial conditions specified above, there results
I1 w ¼ I3 na
Substituting the expressions of fη, fZ, and w` resulting from the six scalar
equations, that is,
fg ¼ mRða þ xz wÞ
fZ ¼ mg
w ¼ I3 na=I1
into the first of these six equations, which is rewritten below for convenience
I1 a I3 xz w ¼ fZ Ra þ fg R
there results

2 2 I3 þ mR
2
I1 þ mR a þ I3 n mgR a ¼ 0
I1
This differential equation shows that the plane of the disc wobbles in and out of
the vertical, provided that the value of the expression within square brackets is
greater than zero, that is, provided that
mgRI1
n2 [
I3 ðI3 þ mR2 Þ
This is because, when the condition written above is satisfied, then it is possible
to define a quantity k2 such that
I3 þ mR2
I 3 n2 mgR
I1
k2 ¼
I1 þ mR2
and put the differential equation

2 I3 þ m R2
I1 þ m R a þ I3 n
2
mgR a ¼ 0
I1
in the following form
a þ k2 a ¼ 0
As is well known, the preceding differential equation has the following solution
a ¼ A cosðktÞ þ B sinðktÞ
where the values of A and B depend on the initial conditions, and k is the angular
frequency of the periodic motion of the plane of the disc.
In addition, the expression found above I1w` = I3na indicates that w` is pro-
portional to a, and therefore, w` also wobbles sinusoidally. In other words, the disc
rolls in a wavy line, which is nearly straight.
(b) Disc spinning about the vertical axis

In this case, the main motion of the disc is a spin about the vertical axis. As a result
of disturbances, the disc moves in a small circle, such that the value of xz is
negligible in comparison with the value of w`. The value of the angle h is, as was in
the case discussed above, nearly equal to p/2, and therefore, the value of the angle
a = p/2 − h is nearly equal to zero.
In these conditions, the second of the six general equations, that is,
MOy ¼ I1 ðw sin h þ wh cos hÞ þ I1 hw cos h þ I3 xz h ¼ 0
becomes
I1 w ¼ 0
which means w` = constant, and the other five equations become
I1 a I3 xz w I1 w2a ¼ fZ Ra þ fg R

I3 xz ¼ fx R
fx ¼ mRðxz awÞ
fg ¼ mRðaþ xz wÞ
fZ mg ¼ 0
The two equations
I3 xz ¼ fx R
fx ¼ mRðxz awÞ
yield
ðI3 þ mR2 Þxz ¼ mR2 aw
By integrating the preceding differential equation with the initial conditions

a = 0 and xz = 0 for t = 0, there results
ðI3 þ mR2 Þxz ¼ mR2 wa
The preceding equation, solved for xz, yields
xz ¼ mR2 wa=ðI3 þ mR2 Þ
Substituting this expression of xz and also fη = mR(a`` + xzw`) and fZ = mg into

the equation −I1 a`` − I3 xz w` − I1 w`2 a = −fZ R a + fη R, there results
ðI1 þ mR2 Þa þ ½ðI1 þ mR2 Þw2 mgRa ¼ 0

which means that, in order for the spinning motion of the disc to be stable, the value
of the quantity between square brackets must be greater than zero, and therefore, the
following inequality must be satisfied:
mgR
w2 [
I1 þ mR2
When this happens, then the disc oscillates sinusoidally with a small angle a
about the x-axis, which in turn spins about the vertical Z-axis at a speed w`.
(c) Disc spinning nearly horizontally

The problem considered below concerns a disc (e.g. a coin) spinning about a
vertical axis, in such a way that the plane of the disc is nearly horizontal. The
frequency of the sound generated by the disc on the floor increases very rapidly
during the last stage of oscillation. In addition, the point of contact with the floor
spins around a circle, whose diameter is nearly equal to the diameter of the disc, and
the component xz of the angular velocity vector x has a very small value; that is,
the face of the disc rotates very slowly. The z-axis of the body-fixed system xyz is
nearly vertical, and therefore, the value of the angle h is very small. The z-axis has a
very rapid motion of precession around the vertical Z-axis of the immovable XYZ
system, and therefore, the value of w` is very large.
The six scalar equations which express the moments and the forces are
I1 h I1 w2 h ¼ fZ R þ fg Rh

I1 wh þ 2I1 wh ¼ 0
I3 xz ¼ fx R
fx ¼ mRxz
fg ¼ mRxz w
fZ mg ¼ mRh
The two equations I3 x`z = fx R and −fx = m R x`z imply

I3 þ mR2 xz ¼ 0
and therefore
xz ¼ constant ¼ n
The equation I1 w`` h + 2 I1 w` h` = 0 may be rewritten as follows
w h
¼ 2
w h
The preceding differential equation, integrated with the initial conditions

w` = w`0 and h = h0 for t = 0, yields
" #
w h h 2
ln ¼ 2 ln ¼ ln
w0 h0 h0

2
h
w ¼ w0
h0
Substituting this expression of w` and also fη = m R xz w` and fZ = mg + m R h``

into the equation I1 h`` − I1 w`2 h = −fZ R + fη R h, there results
I1 w20 h40 mR2 nw0 h20

I1 þ mR2 h ¼ mgR þ
h 3 h
Since h3 and h appear in the denominators on the right-hand side of the pre-
ceding differential equation, then the value of the acceleration h`` increases and
tends to infinity as the value of h approaches zero.
Finally, the equation w` = w`0(h/h0)−2 indicates that the speed of precession w`
also increases as h decreases.
References
1. L. Landau, E. Lifchitz, in Mécanique, Editions Mir, Moscou (1969)

2. W.T. Thomson, Introduction to Space Dynamics (Dover Publications, New York, 1986). ISBN
0-486-65113-4
3. J. Diebel, in Representing Attitude: Euler Angles, Unit Quaternions, and Rotation Vectors.
Stanford University, 35 p, 20 Oct 2006. Web site www.swarthmore.edu/NatSci/mzucker1/e27/
diebel2006attitude.pdf
4. C.W. Martz, in Attitude reorientation of spacecraft by means of impulse coning. NASA
Technical Note D-8452, 35 p, Aug 1977
5. L. Poinsot, in Théorie nouvelle de la rotation des corps, Bachelier, Paris (1852). Web site
https://archive.org/details/thorienouvelled01poingoog
6. A.G. Greenhill, The Applications of Elliptic Functions (Macmillan & Co., London, 1892)
Chapter 8
Instruments for Aerospace Navigation
8.1 Motion of a Symmetric Gyroscope
A mechanical gyroscope is a spinning wheel or disc mounted in a movable sus-

pension frame, which allows the axis of rotation of the wheel to point in any
direction. It looks like a spinning top. The following figure, due to the courtesy of
NASA [1], illustrates a mechanical gyroscope.
Gyroscopes are used on space vehicles for two principal purposes:

• to establish an inertial system of reference; and
• to measure the orientation angles and their rates of change with time of a space
vehicle with respect to an inertial reference system.
A gyroscope is mounted on a gimbal set, which is a series of concentric rings,
each of which can pivot about a single axis, in order to allow the rotation of an

https://doi.org/10.1007/978-3-319-62220-0_8
1042 8 Instruments for Aerospace Navigation
object about that axis, as shown in the following figure, which is also due to the
For example, a two-degree-of-freedom gyroscope has a support which allows the

spin axis of the gyroscope to have two degrees of rotational freedom. For this
purpose, this gyroscope is supported with a gimbal (the outer gimbal) between the
gyroscope element (the inner gimbal) and the case.
The following figure is a scheme of the two-degree-of-freedom gyroscope
illustrated above. Let z be the spin axis of the rotating wheel. The inner gimbal
allows a pitching rotation of the spin axis about the horizontal bearings x, whereas
the outer gimbal is free to rotate about the vertical Z-axis. Then, the gyroscope is
pivoted about the stationary geometric centre of the gimbal system. As shown in the
figure, the centre of mass G of the rotor does not coincide with the fixed centre O.
The distance between these two points is denoted with ‘ in the following figure.
8.1 Motion of a Symmetric Gyroscope 1043
The rotation w` of the horizontal x-axis (node axis) about the vertical Z-axis is
called precession. When the angle h, which the spin axis z forms with the vertical
Z-axis, is held constant, then the spin axis generates a cone because of the motion of
precession. The rotation h` of the inner gimbal about the node axis x is called
nutation. The precession and nutation of the spin axis of a gyroscope are like the
precession and nutation of the rotation axis of the Earth, as will be shown in
Sect. 8.3. In the general case, these two motions may exist simultaneously.
Goldstein et al. [3] point out that the torque-free precession of the axis of the Earth
is not to be confused with its slow precession about the normal to the ecliptic plane.
This astronomical precession of the equinoxes is due to the gravitational torques
exerted by the Sun and Moon. The long period of the precession of the equinoxes
(about 26000 years) is to be compared with a period of roughly one year for the
torque-free precession.
As a first approximation, we neglect the masses of the gimbals with respect to
the mass of the rotating disc. Therefore, the disc is free to rotate in any possible
manner about the fixed centre O, as shown in the preceding figure. In these con-
ditions, a gyroscope is identical to a spinning top pivoting about a fixed point O and
is only subject to a torque MOx = mg‘ sin h acting about the x-axis due to the weight
force mg. In the particular case of ‘ = 0, a gyroscope is free to rotate about its centre
of gravity. It is convenient to write the Euler equations in the reference system xyz
attached to the rotating body, taking the moments with respect to the fixed origin O,
such that the line of nodes (x-axis) is one of the three perpendicular axes. For this
purpose, it is necessary to express the three components of the angular velocity
vector x and the three components of moment-of-momentum vector hO. As shown
in the preceding figure, the angular velocities of the three axes x, y, and z are
respectively
xx ¼ h
xy ¼ wsin h
xz ¼ wcos h
Since the rotating disc is an axially symmetric body, let I1, I1, and I3 be its
moments of inertia along the principal axes of inertia, respectively, 1, 2, and 3.
As also shown by the preceding figure, the three components of the
moment-of-momentum vector hO ¼ = x along the axes x, y, and z of the
body-fixed system are respectively
hOx ¼ I1 h
hOy ¼ I1 wsin h
hOz ¼ I3 ð/þ wcos hÞ
where / is the spin angle of the principal axes of inertia 1 and 2 with respect to the
line of nodes x. We can now form the vector product x hO, which is on the
right-hand side of the fundamental equation of Sect. 7.6, which is written below:
M O ¼ hO þ x hO
The components xx, xy, and xz of the angular velocity vector x have been
expressed above. The vector product on the right-hand side yields
x hO ¼ ðhOz wsin h hOy wcos hÞux þ ðhOx wcos h hOz hÞuy

þ ðhOy h hOx wsin hÞuz
The only torque acting upon the rotating body, as a result of the weight force, is
MOx = mg‘ sin h, because MOy = MOz = 0.
Hence the fundamental equation MO = hÒ + x hO (see Sect. 7.6), projected
onto the axes x, y, and z, yields
MOx ¼ hOx þ hOz wsin h hOy wcos h

MOy ¼ hOy þ hOx wcos h hOz h
MOz ¼ hOz þ hOy h hOx wsin h
Substituting MOx = mg‘ sin h, MOy = MOz = 0, the components of hÒ and their
first derivatives with respect to time into the preceding equations, there results
mg‘ sin h ¼ I1 hþ I3 ð/þ wcos hÞwsin h I1 w2 sin h cos h

0 ¼ I1 ðwsin hÞþ I1 hwcos h I3 hð/þ wcos hÞ
0 ¼ I3 ð/þ wcos hÞ
As a result of the third equation written above, the quantity /` + w` cos h is

constant in time. This quantity is denoted below by n.
On the other hand, the three components of the angular velocity vector x along
the three principal axes of inertia 1, 2, and 3 of the rotating body are respectively
x1 ¼ hcos / þ wsin h sin /

x2 ¼ hsin / þ wsin h cos /
x3 ¼ wcos h þ / ¼ n
By adding the squares of x1 and x2, there results
x21 þ x22 ¼ h2 þ w2 sin2 h
Therefore, the kinetic energy of the rotating body results
1 1 1 1
Trot ¼ I1 x21 þ x22 þ I3 x23 ¼ I1 h2 þ w2 sin2 h þ I3 n2
2 2 2 2
By taking the value zero of the potential energy, U, of the rotating body at the
height of the plane XY, the value of the potential energy at a height H = ‘ cos h
results
U ¼ mg‘ cos h
Therefore, the total energy of the rotating body is
1 1
E ¼ Trot þ U ¼ I1 h2 þ w2 sin2 h þ I3 n2 þ mg‘ cos h
2 2
The preceding equation is one of the first integrals of the differential equations of
motion.
Since the moment MOZ about the Z-axis is equal to zero, then the moment of
momentum hOZ = hOz cos h + hOy sin h must be constant in time. Remembering the
expressions hOz = I3(/` + w` cos h) and hOy = I1w` sin h, there results
hOZ ¼ I3 ð/þ wcos hÞ cos h þ I1 wsin2 h ¼ I3 n cos h þ I1 wsin2 h
Solving the preceding equation for w`, there results
hOZ I3 n cos h
w ¼
I1 sin2 h
which in turn, substituted into the equation of the total energy, yields
1 1 ðhOZ I3 n cos hÞ2

E I3 n2 ¼ I1 h2 þ þ mg‘ cos h
2 2 2I1 sin2 h
The preceding equation is written in terms of h. Let us set for convenience

a = 2(E – 1/2I3n2)/I1 a constant
b = 2mg‘/I1 a constant
c = hOZ/I1 a constant
N = I3n/I1 a constant
u = cos h a function of time
By using these symbols, the equation of the total energy considered above may
be rewritten as follows
a sin2 h ¼ h2 sin2 h þ ðc N cos hÞ2 þ b cos h sin2 h

which, expressed in terms of u, becomes

u2 ¼ ða buÞ 1 u2 ðc NuÞ2
The solution of the equation written above is given by the following integral
Zu
dx
t t0 ¼ h i12
u0 ða bxÞð1 x2 Þ ðc NxÞ2
Fortunately, it is not necessary to evaluate the integral given above to understand

the motion of a gyroscope. For this purpose, by setting u`2 = f(u), the preceding
equation u`2 = (a − bu)(1 − u2) − (c − Nu)2 may be written as follows

f ðuÞ ¼ ða buÞ 1 u2 ðc NuÞ2
As will be shown in the following discussion, the roots of this cubic equation
provide information on the motion of a gyroscope. First of all, in order for a
physical solution to exist, f(u) cannot be negative. The limits of the motion,
expressed in terms of the angle h, are determined by the three roots of the equation
f(u) = 0. Since h is between 0 and p/2, as shown in the preceding figure, then
u = cos h must be between 0 and 1. The value of the largest root, u3, is greater than
1, and the values of the smaller roots, u1 and u2, if they be real, are in the interval
−1 u 1. The function f(u) has positive values only in this interval. The
following figure, due to the courtesy of Professor Richard Fitzpatrick [4], shows a
case in which u1 and u2 are within the interval [0, 1].
For large values of u, the dominant term in the preceding equation is bu3, which
results from the product of the first two terms on the right-hand side of the equation.
Hence, f(u) is positive for large positive values of u, and negative for large negative
values of u. For u = ±1, the quantity (a − bu)(1 − u2) vanishes, and the preceding
equation reduces to
f ð1Þ ¼ ðc NuÞ2
Since f(±1) is equal to a squared quantity with a minus sign in front of it, then
f(u) is negative for u = ±1. On the other hand, since f(u) is by definition
f ðuÞ ¼ u2 ¼ ðcos hÞ2 ¼ h2 sin2 h
then f(u) is positive when h` and h have real values. Consequently, in all practical
cases, the values of the function u = cos h must be within an interval [u1, u2] for
which f(u) is positive, as shown in the preceding figure.
In addition, for h > 0, h` must be equal to zero at u = u1 and u = u2. Therefore,
the tip of the spin axis z of a gyroscope must move between two bounding circles
such that u1 = cos h1 and u2 = cos h2. In other words, h1 and h2 are the limits of the
vertical (in the direction of the Z-axis) motion of a gyroscope. The spin axis of a
gyroscope oscillates backward and forward between h1 and h2 as the gyroscope
precesses about the vertical axis. This oscillation is called nutation.
The type of curve traced out by the tip of the spin axis between the angles h1 and
h2 depends on the values of the constant quantities c and N. These curves are shown
in the following figure, which is due to the courtesy of the University of California
at Los Angeles [5].
In this figure, the central curve of the three nutation curves illustrated above
refers to the case in which c > Nu2. The upper circle and the lower circle indicate,
respectively, the angle h2 and the angle h1 (such that h2 < h1, because u2 > u1). The
angle h is measured from the vertical axis, which represents the spin axis of a
rotating gyroscope. In order to show how the type of curve traced by the tip of the
spin axis depends on the values of the constant quantities c and N, let us rewrite the
preceding equation
hOZ I3 n cos h
w ¼
I1 sin2 h
in terms of the two constants c and N, and of the variable u. This yields
c Nu
w ¼
1 u2
Since the value of u is less than unity between u1 and u2, then the positive or
negative sign of w` depends only on the value of the numerator c – Nu of the
preceding fraction.
In particular, if c = Nu2, then w` is equal to zero when the tip of the spin axis
reaches the upper bounding circle and is positive when h is greater than h2, as
shown by the curve illustrated on the left-hand side of the preceding figure.
The curve illustrated on the right-hand side of the preceding figure refers to the
case in which w` changes sign for some value ui of u such that u1 < ui < u2.
Therefore, this case is described by c − Nui = 0 for u1 < ui < u2.
The following section describes the motion of a gyroscope as a function of the
initial conditions. Let h = h0, and w` = h` = 0 be the conditions at t = 0. In these
conditions, the values of the two constants hOZ and E result by substituting these
initial values into the equations derived above
hOZ I3 n cos h
w ¼
I1 sin2 h
1 1 ðhOZ I3 n cos hÞ2

E I3 n2 ¼ I1 h2 þ þ mg‘ cos h
2 2 2I1 sin2 h
The first equation yields
hOZ ¼ I3 n cos h0
The second equation yields
1
E I3 n2 ¼ mg‘ cos h0
2
Hence, the two equations derived above may be rewritten, in terms of the initial
conditions, as follows
N ðcos h0 cos hÞ
w ¼
sin2 h

N 2 ðcos h0 cos hÞ
h ¼ ðcos h0 cos hÞ b
2
sin2 h
Since the value of the term on the right-hand side of the second equation cannot
be negative, then the angle h0 must correspond to the upper of the two bounding
circles illustrated in the preceding figure, which represent the angle between the
spin axis and the vertical axis.
The value of h corresponding to the lower of the two bounding circles can be
found by setting h` = 0 and solving the second equation for cos h. This yields
" 2 2 #12
N2 N2 N
cos h ¼ 1 cos h0 þ
2b b 2b
where the sign which takes effect in front of the square root must be negative, as
Since cos h0 is less than 1, then the square root in the preceding expression is
greater than
" 2 2 #12 " 2 2 #12
N2 N N2 N2 N N2
1 þ ¼1 \ 1 cos h0 þ ¼ cos h
b 2b 2b b 2b 2b
Therefore, if the positive sign were used in front of the square root, then cos h
would be greater than 1, which is impossible. For this reason, the angle h1, which
corresponds to the lower bounding cycle, is given by
" 2 2 #12
N2 N2 N
cos h1 ¼ 1 cos h0 þ
2b b 2b
8.2 Steady Precession of a Symmetric Gyroscope
As has been shown in Sect. 8.1, a symmetric rotating gyroscope moves so that the
tip of its axis of symmetry, z, is constantly within a zone comprised between the
two angles h1 and h2 (measured from the vertical Z-axis), which correspond to the
two roots u1 and u2 of the cubic equation f(u) = 0.
If the roots u1 and u2 become closer to each other, then the zone between h1 and
h2 narrows, and reduces to zero when u1 = u2, that is, in case of the cubic equation
having three real roots, two of which are coincident. In this case, the curve rep-
resenting the cubic polynomial is tangent to the u-axis in the point of abscissa
u1 = u2, as shown in the following figure.
Let hs be the value of h corresponding to the case h1 = h2, where the subscript
s stands for steady. The angle hs is called the angle of steady precession, and results
from the initial conditions h = hs, h` = 0, and w` = w`s.
In case of steady precession, there results h` = 0, and consequently the three
components of the angular velocity vector x along the three axes x, y, and z of the
body-fixed system are respectively
xx ¼ h ¼ 0
xy ¼ wsin h
xz ¼ wcos h
The three components of the moment-of-momentum vector hO along the axes

x, y, and z are respectively
hOx ¼ I1 h ¼ 0
hOy ¼ I1 wsin h
hOz ¼ I3 ð/þ wcos hÞ ¼ I3 n
where n = /` + w` cos h has been found in Sect. 8.1 to be a quantity constant in

time. Likewise, the following equation derived in Sect. 8.1
becomes in the present case
mg‘ sin h ¼ I3 nwsin h I1 w2 sin h cos h

8.2 Steady Precession of a Symmetric Gyroscope 1051
and may also be written as follows
I3 n mg‘
w2 wþ ¼0
I1 cos h I1 cos h
This equation is represented graphically by a parabola, which is shown below.
The two roots, w`1 and w`2, of the preceding equation are
" 2 #12
I3 n I3 n mg‘
w1;2 ¼
2I1 cos h 2I1 cos h I1 cos h
These two roots represent the two velocities of precession. In order for the
expression under the square root to be positive, the following condition must be
satisfied
4I1 mg‘ cos h

n2 [
I32
The equation shown above and its two roots can also be represented graphically
by a plot of the moment MOx about the x-axis against the velocity of precession w`,
as shown in the preceding figure.
To any positive value of the moment MOx, there correspond two values of the
velocity of precession w`. When the value of MOx is equal to zero, these two values
are 0 and w`0, the latter value being given by
I3 n I3 /
w0 ¼ ¼
I1 cos h ðI3 I1 Þ cos h
which is just the equation found in Sect. 7.8 for an axially symmetric (I1 = I2)
rotating body not subject to external moments.
The velocity of precession corresponding to the maximum value of the moment
MOx about the x-axis is
1 I3 n
wp ¼ w0 ¼
2 2I1 cos h
The maximum value of the moment MOx corresponding to w`p is
1 I32 n2
ðMOx Þmax ¼ mg‘ sin h ¼ tan h
4 I1
As shown above, to any positive value of MOx < (MOx)max there correspond two
velocities of precession, w`1 and w`2, which are called the slow precession and the
fast precession. In practice, the fast precession is not reached because of the high
value of the kinetic energy required. Therefore, the precession of a spinning
gyroscope is usually of the slow type.
As has also been shown above, the following condition must be satisfied
4I1 mg‘ cos h

n2 [
I32
in order for the two roots w`1 and w`2 to have real values. Therefore, to this effect, in
case of h being equal to 90°, n must be greater than zero. However, in this case, the
two roots w`1 and w`2 are indeterminate, because cos h = 0 appears at the
denominator in the expression of w`1,2 given above. The problem of a gyroscope
rotating with an angle h = 90° can be solved by considering the following figure.
According the fundamental equation of the dynamics of rotating bodies, the rate
of change hÒz of the component along z of the moment-of-momentum vector hO is
equal to the moment applied MOz, which is equal to mg‘. Since hOz = I3n, then
8.2 Steady Precession of a Symmetric Gyroscope 1053
hOz ¼ I3 nw ¼ mg‘
Therefore, the velocity w`, in case of h = 90°, is
mg‘
w ¼
I3 n
and this velocity is finite as long as n is finite.

The same result can also be found by considering the equation derived above
I3 n mg‘
w2 wþ ¼0
I1 cos h I1 cos h
When cos h tends to zero because h tends to 90°, then the quadratic term w`2
becomes negligible in comparison with the other two terms of the equation.
Another limiting case (h = 0°) is known as that of a sleeping top, meaning by
this name a top rotating about an axis directed vertically upward.
In this case, the value of n resulting from the general equation derived above
4I1 mg‘ cos h

n2 [
I32
must be such that

1
2ðI1 mg‘Þ2
n
I3
which condition is often used to determine the spin to be given to a missile or

projectile for stability. When the velocity vector of a projectile acts along the axis of
symmetry of the projectile, then the drag force also acts along the axis of symmetry.
By contrast, when the axis of symmetry deviates by an angle h, even small, from
the direction of the velocity vector, then the drag force D (which acts in the centre
of pressure, which is located at a distance ‘ ahead of the centre of mass) generates a
moment MO = D‘ sin h about the centre of mass O.
Therefore, in order for a projectile to be in stable motion along its trajectory, a
spin is applied to the projectile, such that
1
2ðI1 D‘Þ2
n
I3
For this purpose, the barrel of a gun or firearm has helical grooves, which impart
a spin to the projectile about its axis of symmetry. This spin makes the motion of
the projectile stable by gyroscopic effect.
8.3 Precession and Nutation of the Polar Axis of the Earth
The effects induced on the motion of an artificial satellite by the non-spherical shape
of the Earth have been shown in Sect. 3.2 et seq. According the third principle of
dynamics, the same effects are induced, in the opposite direction, by a satellite on
the motion of the Earth. In practice, since any artificial satellite has a mass much
smaller than the mass of the Earth, then the effects induced by this satellite on the
motion of the Earth are negligible. Things change when these effects are induced on
the Earth by natural bodies of the Solar System whose masses are not negligible in
comparison with the mass of the Earth.
The Earth may be considered as a spinning top, whose axis of symmetry has a
motion of precession about an axis perpendicular to the ecliptic plane. This motion,
which is known as the precession of the equinoxes, is due to the non-spherical
shape of the Earth. If the Earth were perfectly spherical, none of the bodies
belonging to the Solar System could induce gravitational torques on it.
The Earth, which is approximately an oblate ellipsoid of revolution, is subject to
gravitational torques which are due primarily to attractive forces exerted by the Sun
and the Moon on the equatorial bulge of the Earth. These torques cause the polar
axis of the Earth to have a motion of precession perpendicularly to the ecliptic
plane. In order to determine the gravitational attraction due to a nearly spherical
body, let us consider the body shown in the following figure.
We want to compute the potential of gravitation in an external point P due to the

attractive force exerted by the total mass M of the body. Let us consider a system of
reference xyz, whose origin O coincides with the centre of mass of the given body.
Let x, y, z and n, η, f be the co-ordinates of, respectively, the point P and the
point-mass dM located within the body.
As has been shown in Sect. 3.2, the potential of gravitation in P due to the
point-mass dM = qdv = qdndηdf is
ZZZ
qðn; g; fÞdndgdf
V ¼G h i12
v
ðx nÞ2 þ ðy gÞ2 þ ðz fÞ2
8.3 Precession and Nutation of the Polar Axis of the Earth 1055
(Heiskanen and Moritz [6]), where q is the density and v is the volume of the body
considered. Let
1
R ¼ x2 þ y2 þ z 2 2
1
s ¼ n2 þ g2 þ f 2 2
h i12 1
r ¼ ðx nÞ2 þ ðy gÞ2 þ ðz fÞ2 ¼ R2 þ s2 2Rs cos c 2
be the lengths of the segments shown in the preceding figure, where c is the angle
between the two position vectors of, respectively, the point-mass dM and P. Hence
s 2 s 12
1 1 1
¼ ¼ 1þ 2 cos c
r ðR2 þ s2 2Rs cos cÞ12 R R R
As has been shown in Sect. 3.2, by expanding 1/r in a series of Legendre

polynomials, there results
s 2 s 12 1 i
1 1 1X s
¼ 1þ 2 cos c ¼ Pi ðcos cÞ
r R R R R i¼0 R
where Pi(x) is the Legendre polynomial of ith degree. This series converges when
R > s. This is because, after setting for convenience s/R = x, there results
1 1 1
¼ ð1 zÞ2
1
1 ¼ 1 ¼ 1
ð1 þ x2 2x cos cÞ 2 ½1 xð2 cos c xÞ 2 ð1 zÞ 2
where z = x(2 cos c − x) = (s/R)(2 cos c − s/R). Now, we expand f(z) = (1 − z)−1/2
in a Maclaurin series around z0 = 0. To this end, we evaluate the function f(z) and
its derivatives up to the fourth order at z = 0. Using the notation df(z)/dz
z′, we
have
h i h i
½f ðzÞ 0 ¼ ð1 zÞ2 ¼ 1 ½f 0 ðzÞ 0 ¼ 12 ð1 zÞ2 ¼ 12
1 3
h 0i

ð3Þ h15 0 i
½f 00 ðzÞ 0 ¼ 34 ð1 zÞ2 ¼ 34 f ðzÞ 0 ¼ 8 ð1 zÞ2 ¼ 15
5 7
8

ð4Þ h105 0 i
92
0
f ðzÞ 0 ¼ 16 ð1 zÞ ¼ 16
105
0
Then the series expansion of (1 − z)−1/2 around z0 = 0 is

1 1 3 2 1 15 3 1 105 4
ð1 zÞ2 ¼ 1 þ
1
zþ z þ z þ z þ
2 2! 4 3! 8 4! 16
1 3 2 5 3 35 4
¼ 1þ zþ z þ z þ z þ
2 8 16 128
Hence
1 3
½1 xð2 cos c xÞ 2 ¼ 1 þ xð2 cos c xÞ þ x2 ð2 cos c xÞ2
1
2 8
5 3 3 35 4
þ x ð2 cos c xÞ þ x ð2 cos c xÞ4 þ
16 128
Expanding the binomial (2 cos c − x) in a series of powers yields
1 3
½1 xð2 cos c xÞ 2 ¼ 1 þ x cos c x2 þ x2 4 cos2 c 4x cos c þ x2
1
2 8
5 3
þ x 8 cos3 c 12x cos2 c þ 6x2 cos c x3
16
35 4
þ x ð16 cos4 c 32x cos3 c þ 24x2 cos2 c
128
8x3 cos c þ x4 Þ þ
The expression written above, ordered by increasing powers of x, becomes
1 1
½1 xð2 cos c xÞ 2 ¼ 1 þ x cos c þ x2
1
3 cos2 c 1 þ x3 ð5 cos3 c
2 2
1
3 cos cÞ þ x4 35 cos4 c 30 cos2 c þ 3 þ
8
½1 xð2 cos c xÞ 2 ¼ x0 P0 ðcos cÞ þ x1 P1 ðcos cÞ þ x2 P2 ðcos cÞ

1
þ x3 P3 ðcos cÞ þ x4 P4 ðcos cÞ þ
where the nth expression which follows the nth power of x, denoted by Pn(cos c), is
the Legendre polynomial of the nth order.
When R is much greater than s, then the series expansion can be truncated after
the second power of x = s/R. By so doing, there results
s 2 s 12 s s 2 1
1 1
1þ 2 cos c ¼ 1þ cos c þ 3 cos2 c 1
R R R R R R 2
s 3
þO
R
In this case, the potential of gravitation may be approximated as follows

Z s s 3
G 1 s 2 3 s 2 2
V¼ 1þ cos c þ cos c þ O dM
R R 2 R 2 R R
M
The first four addends

R under the sign of integral
R give rise to the following
integrals: (1) ðG=RÞ M dM ¼ GM=R; (2) ðG=R Þ M ðs cos cÞdM ¼ 0, because the
2
axis OP shown in theRpreceding figure passes throughR the centre of mass of the
body; (3) 1=2ðG=R3 Þ M s2 dM; and (4) 3=2ðG=R3 Þ M ðs2 cos2 cÞdM.
Since cos2 c = 1 − sin2c is a trigonometric identity and the second integral is
equal to zero, then the sum of the four integrals written above can be written as
follows
Z Z
GM 1 G 3G 2 2
þ 3
2s 2
dM 3
s sin c dM
R 2R 2R
M M
Since the term 2s2 may be written as follows

1 1
1
2s2 ¼ 2 x2 þ y2 þ z2 2 ¼ 2 x21 þ x22 þ x23 2 ¼ x22 þ x23 þ x21 þ x23 þ x21 þ x22 2
then there results

Z
1G 1G
2s2 dM ¼ ðI1 þ I2 þ I3 Þ
2 R3 2 R3
M
where I1, I2, and I3 are the moments of inertia of the given body about the three
principal axes of inertia, respectively, x1, x2, and x3, as has been shown in Sect. 8.2.
On the other hand, the preceding figure shows that
Z
2 2
s sin c dM ¼ IOP
M
where IOP is the moment of inertia of the given body about the axis OP.
Therefore, the potential of gravitation of a spheroid may be approximated as
follows
GM 1 G s 3
V¼ þ ð I1 þ I 2 þ I 3 3I OP Þ þ O
R 2 R3 R
The preceding expression of the potential of gravitation of a spheroid is known

as MacCullagh’s formula.
When the given body is a sphere, there results I1 = I2 = I3 = IOP, and therefore
the potential of gravitation of a spherical body of mass M is V = GM/R at any point
P. In case of points P whose distance R from the centre of mass O of the given body
is much greater than the distance s of the point-mass dM from O, the potential of
gravitation is V = GM/R.
The Earth and the other planets of the Solar System may be considered as axially
symmetric ellipsoids of revolution, such that two (e.g. I1 and I2) of the three
moments of inertia about the principal axes x1, x2, and x3 are equal (I1 = I2). For
such bodies, as shown in the following figure, the angle c is the geocentric latitude
of the axis OP.
In this case, there results
I1 ¼ I2
IOP ¼ I1 cos2 c þ I3 sin2 c
By using the trigonometric identity cos2 c = 1 − sin2c, the preceding equation

becomes
IOP ¼ I1 ð1 sin2 cÞ þ I3 sin2 c ¼ I1 þ ðI3 I1 Þ sin2 c
and consequently
I1 þ I2 þ I3 3IOP ¼ 2I1 þ I3 3½I1 þ ðI3 I1 Þ sin2 c ¼ ðI3 I1 Þð1 3 sin2 cÞ
Therefore, MacCullagh’s formula for an axially symmetric rotating body is
GM 1 G s 3
V¼ 3 ðI3 I1 Þ 3 sin2 c 1 þ O
R 2R R
This formula can be compared with the following second-order formula of

Sect. 3.2.
r 2
l 1
C20 3 sin2 u 1
E
V¼ 1þ
R r 2
where l = GM is the gravitational parameter of the Earth, G is the universal

gravitational constant, r is the radius vector from the centre of mass of the Earth to
the point P, rE and M are, respectively, the equatorial radius and the mass of the
Earth, C20 = −0.001082636 is a constant, and u* is the geocentric latitude of the

point P where the potential of gravitation is required. This comparison shows that
I3 I1
C 20 ¼
MrE2
The expression 3 sin2u* − 1 is the Legendre polynomial of the second degree.

Let us consider again MacCullagh’s formula for an axially symmetric rotating
body. The only term of the potential of gravitation depending on the orientation of
the body and therefore capable of generating torques is
1G
ðI3 I1 Þ 3 sin2 c 1
2 R3
As the Sun, the Moon, and any other attracting body change their positions with
respect to the Earth, the angle c also changes. In practice, the gravitational torques
due to the Sun and the Moon contribute in the greatest amount to the precession of
the Earth. This phenomenon is illustrated in the following figure, which shows how
the attractive force due to an external body generates a torque which affects the axis
of rotation of the Earth, as a result of the equatorial bulge of the latter. The direction
of this motion of precession is retrograde, that is, clockwise if viewed from a point
placed above the plane of the ecliptic.
The big orbit shown in the preceding figure is the path of the Sun in its apparent
motion with respect to the Earth in the ecliptic plane. The same physical phe-
nomenon is due to the Moon, whose orbit is inclined nearly 5° to the ecliptic plane.
According to the second principle of dynamics for rotating bodies, the rate of
change of the moment of momentum with time is equal to the external torque due to
the attractive force of the Sun or Moon. The angle e of semi-aperture of the
precession cone is the obliquity of the ecliptic with respect to the equator of the
Earth. The motion of precession is governed by the following equations. Following
Williams [7], let m and r be, respectively, the mass and the distance of the attracting
body from the centre of mass of the Earth. Let M and I1, I1, and I3 (with I1 < I3) be,
respectively, the mass and the moments of inertia of the Earth with respect to the
principal axes. Let us consider a system of reference whose origin is the centre of
mass of the Earth, whose z-axis is aligned with the principal axis of the Earth
corresponding to the maximum moment of inertia I3, whose x-axis lies on the
equatorial plane along the intersection between the ecliptic plane and the equatorial
plane, towards the dynamical vernal point, and whose y-axis also lies on the
equatorial plane 90° east of the x-axis and points towards the summer solstice.
Hence, x, y, and z are the three components of the position vector r, whose origin is
the centre of mass of the Earth and whose tip is the centre of mass of the external
body. The potential energy of the external body in the gravity field of the oblate
Earth is

M 1 1
V ¼ Gm ðI3 I1 Þ 3 sin2 d 1 3
r 2 r
where d and a are, respectively, the declination and the right ascension of the
external body. A simple inspection shows that the preceding expression of V comes
from MacCullagh’s formula derived above. Williams [7] expresses the torque MO,
which acts on the oblate Earth because of the gravitational force due to external
body, as follows
2 3 2 3
sin a yz
3GmðI3 I1 Þ sin d cos d 4 ð I I Þ
cos a 5 ¼ 3Gm 4 xz 5
3 1
M O ¼ r $V ¼
r3 r5
0 0
where the differential operator nabla ($) indicates the gradient of V.

The rate of change of the moment of momentum with time is equal to the torque
whose expression is given above. When the path of the external body with respect
to the Earth is known, then the resulting effects (precession and nutation) on the
Earth can be computed. The positions of the Sun, the Moon, and the other planet are
known with respect to the ecliptic plane, as has been shown in Sects. 5.10 and 5.11.
Let X, Y, Z, and e be, respectively, the geocentric ecliptic co-ordinates of the
external body and the obliquity of the ecliptic at the epoch of interest. The con-
version from X, Y, Z to x, y, z is performed as follows
2 3 2 3
x X
4 y 5 ¼ 4 Y cos Z sin 5
z Z cos þ Y sin
In the torque vector MO, the products yz and −xz of the geocentric equatorial
co-ordinates x, y, and z yield
2 3 2 3
yz 1=2ðY 2 Z 2 Þ sinð2Þ þ YZ cosð2Þ
4 xz 5 ¼ 4 XZ cos XY sin 5
0 0
The geocentric ecliptic co-ordinates X, Y, and Z of the external body can be

expressed as a function of the geocentric distance r and the geocentric ecliptic
longitude k and latitude b, as follows
2 3 2 3
X cos b cos k
4 Y 5 ¼ r 4 cos b sin k 5
Z sin b
Since the path of the Earth about the Sun is an ellipse in the ecliptic plane, then
the gravitational torque due to the Sun can be computed accurately.
The average component MOx of the torque vector MO about the x-axis, over an
integral number of revolution, is
3GmðI3 I1 Þ sin cos

MOx ¼ 3
2a3 ð1 e2 Þ2
where a and e are, respectively, the major semi-axis and the eccentricity of the orbit
of the Earth about the Sun, and e is the obliquity of the ecliptic.
The rate dw/dt of the retrograde precession along the ecliptic is a function of
MOx, as follows
dw MOx 3GmðI3 I1 Þ cos

¼ ¼
dt I3 xz sin 3
2a3 ð1 e2 Þ2 I3 xz
where xz is the principal component of the angular velocity vector x of the Earth, and
I3xz approximates the total moment of momentum relating to the spin of the Earth.
The corresponding expression relating to the precession due to the gravitational
torque generated by the Moon contains an inclination factor of (1 32 sin2 i). The
other two components, MOy and MOz, of the torque vector MO have zero average,
but the first two components have time variations which cause periodic nutation
terms. The approximate model described above works well for the precession
induced by the Sun of the equator of the Earth along the ecliptic, but is a coarse
approximation for the precession induced by the Moon, because the orbit of the
Moon is strongly perturbed by the Sun. Without going into further particulars, we
summarise the preceding discussion as follows. Torques acting on the oblate Earth
due to the gravitational attractions of the Sun and Moon cause motions of pre-
cession and nutation on the equator (or, which is the same, on the spin axis) of the
Earth with respect to extremely distant objects, such as quasars. As a result of these
torques, the tip of the spin axis of the Earth appears to trace a circle in the celestial
sphere over a long period of time (approximately 26,000 years), as shown in the
following figure, due to the courtesy of NASA [8]. The precession is retrograde
with a rate of about 50 seconds of arc per year. Roughly speaking, 1=3 of it is due to
the Sun and 2=3 to the Moon. The rate depends on the masses and distances of the
Sun and Moon, the orbital eccentricities and inclinations, and the obliquity of the
ecliptic with respect to the equator [7].
In addition to the motion of precession, there are small periodic oscillations
superimposed to the conic trajectory described by the spin axis of the Earth, which
are called nutations. They are also shown in the following figure.
The nutations occur because the Moon revolves about the Earth in a tilted
(nearly 5°.14 with respect to the ecliptic plane) elliptic orbit. The intersection of the
orbital plane of the Moon with the ecliptic plane is the line of nodes of the Moon.
This line regresses (with respect to the sense of rotation of the Earth and revolution
of the Moon around the Earth), taking on average 18.6134 years to complete one
cycle (nodal cycle of the Moon). In addition, the longitude of the lunar perigee
progresses (that is, revolves eastward) in the orbital plane of the Moon, taking on
average 8.849 years to complete one cycle (apsidal cycle of the Moon). The nodal
cycle of the Moon causes small variations in the lunar torque as the perigee and
apogee come into the plane of the axial tilt of the Earth. These variations are
sufficient to make the Earth rotation axis nod up and down a little. In addition, there
are annual and semi-annual nutation cycles (caused by variations of the solar couple
at perihelion and aphelion), and 13.6-day cycles (caused by variations of the solar
couple at perigee and apogee), as has been shown by Huggett [9].
8.4 Small Oscillations of Gyroscopes
Let us consider again the symmetric (I1 = I2) gyroscope illustrated in Sect. 8.1, and
shown below for convenience of the reader.
A small action, which perturbs the steady state of this gyroscope, induces har-
monic oscillations with respect to the steady state, as the sequel will show.
8.4 Small Oscillations of Gyroscopes 1063
Let us consider a body-fixed system xyz, whose origin is a fixed point O, and
whose z-axis coincides with the spin axis of the rotating wheel.
Let XYZ be another system of immovable axes, whose origin O is the same as that of
the system xyz attached to the rotor. The centre of mass G of the rotor does not coincide
with the fixed centre O. The distance between these two points is denoted with ‘.
The three scalar equations, which express the components MOx, MOy, and MOz of
the moment vector MO, are those shown in Sect. 8.1, that is,

0 ¼ I1 ðwsin hÞþ I1 hwcos h I3 hð/þ wcos hÞ
0 ¼ I3 ð/þ wcos hÞ
As a result of the third of these equations, the quantity /` + w` cos h is constant

in time. This constant quantity is denoted below by n, and therefore the three scalar
equations may be written as follows
mg‘ sin h ¼ I1 hþ I3 nwsin h I1 w2 sin h cos h

0 ¼ I1 wsin h þ 2I1 hwcos h I3 nh
0 ¼ I3 n
Following Thomson [10], let h0 and w`0 be the values of, respectively, h and w`
in the steady state; and let h* and w`* be the deviations of, respectively, h and w`
from the values h0 and w`0 relating to the steady state. In other words, the
instantaneous values of h and w` are expressed as follows
h ¼ h0 þ h
w ¼ w0 þ w
In case of small oscillations, the following approximations can be made:
hw ¼ h ðw0 þ w Þ h w0

sin h sin h0 þ h cos h0
cos h cos h0 h sin h0
Substituting these expressions into the second of the three scalar equations
which is rewritten below
0 ¼ I1 wsin h þ 2I1 hwcos h I3 nh
and neglecting the products and the squares of the small deviations h* and w`*,
this equation becomes
0 ¼ I1 w sin h0 þ 2I1 h w0 cos h0 I3 nh
By integrating the preceding differential equation, there results
Z
w Zh
I1 sin h0 dw ¼ ðI3 n 2I1 w0 cos h0 Þ dh
0 0
hence
I1 w sin h0 ¼ ðI3 n 2I1 w0 cos h0 Þh
By operating the same substitution into the first of the three scalar equations
which is rewritten below for convenience
mg‘ sin h ¼ I1 hþ I3 nwsin h I1 w2 sin h cos h
and neglecting the products and the squares of the small deviations h* and w`*,
this equation becomes
mg‘ðsin h0 þ h cos h0 Þ ¼ I1 h þ I3 nðw0 þ w Þðsin h0 þ h cos h0 Þ

I1 w0 2 þ 2w0 w sin h0 cos h0 h sin2 h0 þ h cos2 h0
8.4 Small Oscillations of Gyroscopes 1065
For steady precession, h`` is zero, and the preceding equation, written for the
steady state, is
mg‘ ¼ I3 nw0 I1 w0 2 cos h0
Cancelling these terms from the equation
mg‘ðsin h0 þ h cos h0 Þ ¼ I1 h þ I3 nðw0 þ w Þðsin h0 þ h cos h0 Þ

I1 w0 2 þ 2w0 w sin h0 cos h0 h sin2 h0 þ h cos2 h0
and taking account of the equation
solved for w`*, there results

I12 h þ ½ðI3 nÞ2 4I1 mg‘ cos h0 þ I12 w0 2 1 cos2 h0 h ¼ 0
which is a differential equation of the second order of the same type of the
well-known equation
xþ k2 x ¼ 0
Therefore, the oscillations (nutations) h of the spin axis of a symmetric gyro-

scope are sinusoidal and have the following period
2p 2pI1
T0 ¼ ¼h i12
k
ðI3 nÞ2 4I1 mg‘ cos h0 þ I12 w20 ð1 cos2 h0 Þ
When n has so large values that (I3n)2 −4I1mg‘ cos h0 + I21w`20(1 − cos2h0),
then the period of nutation T0 resulting from the preceding expression may be
2pI1
T0
I3 n
The motion of precession has the same period as that of the motion of nutation,
because the equation
solved for w`* yields


I3 n 2I1 w0 cos h0
w ¼ h
I1 sin h0
and therefore w`* is proportional to h*.

When cos h0 = 1, then a gyroscope behaves like a sleeping top (see Sect. 8.2). In
order for the term (I3n)2 − 4I1mg‘ cos h0 + I21w`20 (1 − cos2h0), which appears under
square root in the expression of the period
2pI1
T0 ¼ h i12
2
ðI3 nÞ 4I1 mg‘ cos h0 þ I12 w20 ð1 cos2 h0 Þ
to be greater than or equal to zero, (I3n)2 must be at least equal to 4I1mg‘.

This result confirms the condition
1
2ðI1 mg‘Þ2
n[
I3
to be satisfied for the stability of a sleeping top, as has been found in Sect. 8.2.
8.5 Oscillations of Gyroscopes About Gimbal Axes
Let us consider the two-gimbal gyroscope shown in the following figure. As shown
in the figure, the gimbal axes are x and Z. The centre of mass, O, of the rotor
coincides with the geometric centre of the gimbals.
Neglecting the masses of the gimbals, let I1, I1, and I3 be the moments of inertia
about the axes, respectively, x, y, and z of the system of reference attached to the
rotor and having its origin in the centre of mass O.
Since the scalar angular velocities xx, xy, and xz and moments of momentum
hOx, hOy, and hOz are the same as those shown in Sect. 8.1, their expressions are
rewritten below
MOx ¼ I1 hþ I3 ð/þ wcos hÞwsin h I1 w2 sin h cos h

MOy ¼ I1 ðwsin hÞþ I1 hwcos h I3 hð/þ wcos hÞ
MOz ¼ I3 ð/þ wcos hÞ
8.5 Oscillations of Gyroscopes About Gimbal Axes 1067
With reference to the preceding figure, let us consider the moment MOx about the
horizontal x-axis, and the moments MOy and MOz projected onto the Z-axis.
The projection of MOy and MOz onto the Z-axis yields
MOZ ¼ MOy sin h þ MOz cos h ¼ MOy sin h
This is because MOz = 0. By substituting the expression of MOy written above

into the preceding equation, there results
MOZ ¼ I1 wsin2 h þ 2I1 hwsin h cos h I3 hn sin h
where n = /` + w` cos h. The preceding nonlinear equation may be put in linear

form by making some simplifying assumptions. Since usually the spin velocity /`
is much greater than h` and w`, then the constant n can be approximated to /`. In
addition, the squares and the products of h` and w` can be neglected.
With these simplifying assumptions, the expressions of MOx and MOZ can be
rewritten as follows
MOx ¼ I1 hþ I3 nwsin h

MOZ ¼ I1 wsin2 h I3 hn sin h
There are several solutions of the preceding differential equations depending on

the type of excitation given to the gyroscope.
In case of a steady precession (indicated by the subscript 0), there results h = h0,
and w` = w`0. The quantities h`, h``, and w`` are, all of them, equal to zero. Hence,
for a steady precession, the preceding equations are
MOx0 ¼ I3 nw0 sin h0

MOZ0 ¼ 0
Therefore, for a steady precession, a constant moment MOx0 = I3nw`0 sin h0 is

needed about the horizontal gimbal x-axis.
Let us consider now the case in which a gyroscope, whose spin axis is in steady
precession, is perturbed by a moment which varies as a function of time, so that the
total moment MOx comprises a constant part MOx0 and a variable part MOx(t).
As shown in Sect. 8.1, so in the present case let h* and w`* be the deviations
of, respectively, h and w` about the values h0 and w`0 relating to the steady state.
We express the instantaneous values of h and w` as follows
h ¼ h0 þ h
w ¼ w0 þ w
and make the following approximations
hw ¼ h ðw0 þ w Þ h w0

By so doing, the preceding equations
MOx ¼ I1 hþ I3 nwsin h

MOZ ¼ I1 wsin2 h I3 hn sin h
may be rewritten as follows
h þ I3 nðw0 h cos h0 þ w0 sin h0 þ w sin h0 Þ=I1 ¼ ½MOx0 þ MOx ðtÞ =I1
w sin h0 I3 h n=I1 ¼ 0
The preceding differential equations can be simplified by eliminating the terms

relating to the steady precession, which are defined by
MOx0 ¼ I3 nw0 sin h0

MOZ0 ¼ 0
After cancelling these terms and setting for convenience p = I3n/I1, there results
h þ pðw0 h cos h0 þ w sin h0 Þ ¼ ½MOx ðtÞ =I1

w sin h0 ph ¼ 0
Thomson [10] obtains the solution of the differential equations written above by
means of Laplace transforms, where h* and w`* are the dependent variables.
Following Thomson [11], let f(t) be a known function of t for values of t > 0.
The Laplace transform f ðsÞ
Lf ðtÞ of the function f(t) is defined as follows
Z1
f ðsÞ ¼ est f ðtÞdt
0
where the function f(t) of the real variable t is transformed to a new function f(s) of
the subsidiary variable s, which may be real or complex.
The initial values are h = h0, h` = 0, w` = w`0, h*(0) = h`*(0) = w`*(0) = 0.
The equations involving the transforms are

s2 þ pw0 cos h0 h ðsÞ þ ðp sin h0 Þw ðsÞ ¼ M Ox ðsÞ =I1
ph ðsÞ þ ðsin h0 Þw ðsÞ ¼ 0
This is a system of two algebraic equations for the two unknowns h ðsÞ and
w ðsÞ. By applying Cramer’s rule, the solution of this system is
Ah Aw
h ðsÞ ¼ w ðsÞ ¼
A A
where A, Ah, and Aw are the following 2 2 matrices

ðs2 þ pw0 cos h0 Þ p sin h0
A¼
p sin h0
"
#
M Ox ðsÞ =I1 p sin h0
Ah ¼
0 sin h0
"
#
ðs2 þ pw0 cos h0 Þ M Ox ðsÞ =I1
Aw ¼
p 0
The system of algebraic equations written above can be solved for any given
excitation function MOx(t) and any initial angle h0, which the spin axis z of the rotor
forms with the vertical Z-axis.
When h ðsÞ and w ðsÞ have been determined as has been shown above, the
functions of interest in the domain of time are given by the inverse transforms
h ðtÞ ¼ L1 h ðsÞ and w ðtÞ ¼ L1 w ðsÞ.
For example, let the initial angle of the spin axis (z) be h0 = p/2. Let us suppose
that an impulsive moment M Ox ðtÞ ¼ MdðtÞ is applied to the inner gimbal axis (x),
where d(t) is the Dirac delta function, which is zero everywhere except at t = 0,
where it has an extremely high and extremely narrow peak, that is, its integral
Z
dðtÞdt
is equal to unity over the entire line of time.

Thomson [10] has shown that, in this case, the system of algebraic equations
involving the Laplace transforms has the following solution
M pM
h ðsÞ ¼ w ðsÞ ¼
I1 ðs2
þ p2 Þ I 1 ð s 2 þ p2 Þ
The corresponding functions in the domain of time are

M M
h ðt Þ ¼ sin pt h ðtÞ ¼ cos pt
I1 p I1
M M
w ðtÞ ¼ sin pt wðtÞ ¼ ð1 cos ptÞ
I1 I1 p
Therefore, the actual position of the spin axis of the rotor as a function of time is
p M
þ
hð t Þ ¼ h0 þ h ð t Þ ¼ sin pt
2 I1 p
M
wðtÞ ¼ w0 þ w ðtÞ ¼ ð1 cos ptÞ
I1 p
These results can be interpreted geometrically as follows. With reference to the fol-
lowing figure, let us assume that the spin axis (z) of the rotor is stationary, so that w`0 = 0.
The impulsive moment MdðtÞ shifts suddenly the moment-of-momentum vector

h (whose constant magnitude is h = I3n) along the equator of the sphere (locus of
the tip of h) shown in the preceding figure by an angle
M M
¼
I1 p I3 n
The spin axis of the rotor cannot change instantaneously. It moves downward
from its initial equatorial position at a velocity hð0Þ ¼ M=I1 . The rotation of the
spin axis about the new moment-of-momentum vector h generates a cone, whose
base is a circumference of radius M.
As is well known, in the absence of external moments (in this case, the moment
differs from zero only at t = 0 and is equal to zero at any other time), the
moment-of-momentum vector h is constant and stationary, and the spin axis of the
rotor has a motion of precession around this vector.
Now we consider an initial steady precession of the spin axis of the rotor, such
that the value of the angle h at t = 0 is h0 = p/2 because of an initial moment Mn0.
In this case, the tip of the moment-of-momentum vector h is not along but above
the equator, due to the component I1w`0 of the moment. Again, the impulsive
moment MdðtÞ shifts suddenly the resultant vector h by an amount M along a
circumference whose centre is placed at the latitude
Ix w0
u¼
I3 n
above the equator. At the time following immediately t = 0, the spin axis of the
rotor has an angular velocity whose components are h ð0Þ ¼ M=I1 downward
and w`0 from left to right, as shown in the following figure. Their resultant is
perpendicular to the radial line from the vector h.
The moment-of-momentum vector h, whose magnitude h is approximately I3n,

has a motion of precession along the line of latitude u = Ixw`0/(Izn), at a velocity
w`0. The spin axis rotates about the direction of h along a cone, whose base is a
h i1=2
circumference of radius M2 þ ðI1 w0 Þ2 . The resulting motion is a combination
of nutation and precession. The type of curve described by the tip of the spin axis
depends on the relative values of the two components M=I1 and w`0 of the initial
angular velocity.
8.6 Effects Due to the Moments of Inertia of the Gimbals
Let us consider again the figure of Sect. 8.5, which is also shown below. Due to the
axial symmetry (Ix = Iy) of the rotor, its moments of inertia about the axes x, y, and
z are, respectively, Ix, Ix, and Iz. In addition to these, we take account of the
moments of inertia of the gimbals as follows. The moments of inertia of the inner
gimbal about the axes x, y, and z are, respectively, Ixi, Iyi, and Izi. The moment of
inertia of the outer gimbal about the Z-axis is IZo.
As shown in the preceding figure, the axes x, y, and z rotate at angular velocities
xx = h`, xy = w` sin h, and xz = w` cos h with respect to an immovable reference
system XYZ.
Taking account of the masses which rotate at these angular velocities, the
moments of inertia about the three axes x, y, and z are respectively
8.6 Effects Due to the Moments of Inertia of the Gimbals 1073
I x ¼ Ix þ Ixi
I y ¼ Ix þ Iyi
I z ¼ Izi
The preceding figure also shows that the direction cosines of the Z-axis with
respect to the axes x, y, and z are, respectively, ‘Zx = 0, ‘Zy = sin h, and ‘Zz = cos h.
Therefore, the moment of inertia about the Z-axis is
IZ ¼ IZo þ I x ‘2Zx þ I y ‘2Zy þ I z ‘2Zz ¼ IZo þ I y sin2 h þ I z cos2 h

¼ IZo þ Ix þ Iyi sin2 h þ Izi cos2 h
The components along x, y, and z of the moment-of-momentum vector hO are
hOx ¼ ðIx þ Ixi Þh

hOy ¼ Ix þ Iyi wsin h
hOz ¼ Iz ð/þ wcos hÞ þ Izi wcos h
Remembering the fundamental equation
M O ¼ h0 O ¼ hO þ x hO
of Sect. 1.6, the component
h0 Oz ¼ hOz þ xx hOy xy hOx ¼ hOz þ hhOy ðwsin hÞhOx
of h′O (which is the derivative with respect to time of the moment-of-momentum

vector hO in an immovable reference system XYZ) can be separated in two parts,
which are denoted with [h′Oz]1 and [h′Oz]2, as follows

h0 Oz ¼ ½Iz ð/þ wcos hÞ þ Izi wcos h þ Ix þ Iyi whsin h ðwsin hÞðIx þ Ixi Þh

¼ Iz ð/þ wcos hÞþ Izi ðwcos hÞþ Ix þ Iyi whsin h Ix hwsin h

Ixi hwsin h ¼ ½Iz ð/þ wcos hÞ þ ½Izi ðwcos hÞþ Iyi Ixi hwsin h
¼ ½h0 Oz 1 þ ½h0 Oz 2
In the preceding equation, [h′Oz]1 = Iz(/` + w` cos h)` is the part of h′Oz due to
the axis of the rotor, and [h′Oz]2 = Izi(w` cos h)` + (Iyi − Ixi)h`w` sin h is the part of
h′Oz due to the forces exerted by the outer gimbal on the inner gimbal, as shown in
We assume
½h0 Oz 1 ¼ 0
and consequently
Iz ð/þ wcos hÞ ¼ Iz n ¼ constant
The components hOx, hOy, and hOz of the moment-of-momentum vector hO with
respect to the axes x, y, and z can be projected onto another set of orthogonal axes
x*, y*, and z*, as shown in the following figure.
This projection yields
hOx ¼ hOx ¼ ðIx þ Ixi Þh

hOy ¼ hOy cos h hOz sin h ¼ Ix þ Iyi wsin h cos h ðIz n þ Izi wcos hÞ sin h

hOz ¼ hOy sin h þ hOz cos h þ IZo w ¼ Ix þ Iyi wsin2 h þ ðIz n þ Izi wcos hÞ cos h
þ IZo w
The two equations which express the moments about the axes z* and x* can be
written as follows

Mz
MZ ¼ ðhOz Þ ¼ Ix þ Iyi wsin2 h þ 2whsin h cos h ðIz n þ Izi wcos hÞ
hsin h þ ðIzi wcos h Izi whsin hÞ cos h þ IZo w

¼ Ix þ Iyi sin2 h þ Izi cos2 h þ IZo wþ 2 Ix þ Iyi Izi whsin h cos h
Iz nhsin h ¼ IZ wþ wdðIZ Þ=dt Iz nhsin h ¼ dðIZ wÞ=dt Iz nhsin h

Mx ¼ ðhOx Þ hOy w ¼ ðIx þ Ixi Þhþ Izi Ix Iyi w2 sin h cos h þ Iz wn sin h
¼0
Let us consider the case in which the axes are at rest at the initial time t = 0, and
an initial angular velocity h`0 is given to the z-axis of the inner gimbal, which forms
the initial angle h0 with the Z-axis of the immovable system of reference XYZ. This
initial angular velocity h`0 is given by means of an impulsive moment in the form of
a delta function about the x-axis. The time t = 0 is the instant just after the appli-
cation of the impulsive moment, and therefore Mx = MZ = 0.
The initial conditions are w = 0, w` = 0, h = h0, and h` = h`0 at t = 0. The
assumption w` = 0 at t = 0 is justified, because the initial angular velocity h′0
generates a gyroscopic moment about the Z-axis, as a result of the reaction of the
bearings placed along the x-axis, which moment is not impulsive.
In case of no restraints existing on the Z-axis of the outer gimbal, MZ
Mz* = 0,
and therefore the preceding equation
dðIZ wÞ
Mz ¼ Iz nhsin h
dt

dðIZ wÞ ¼ Iz n sin hdh
The preceding equation canRbe integrated, noting that the value of IZ w` at t = 0

is zero because w` = 0. Since sin hdh ¼ cos h þ c, then this integration yields
IZ w ¼ Iz nðcos h cos h0 Þ

As has been done in Sect. 8.4, so here we set
h ¼ h0 þ h
where h* (which is the deviation of h from h0) is supposed to be small.

Therefore, the following approximations can be made

sin h cos h sin h0 cos h0 þ h ðcos2 h0 sin2 h0 Þ
By taking account of these approximations and setting

I0 ¼ IZo þ Ix þ Iyi sin2 h0 þ Izi cos2 h0
the preceding equations rewritten below for convenience

IZ ¼ IZo þ Ix þ Iyi sin2 h þ Izi cos2 h
IZ wþ Iz nðcos h cos h0 Þ ¼ 0

ðIx þ Ixi Þhþ Izi Ix Iyi w2 sin h cos h þ Iz wn sin h ¼ 0
can be put in the following form

IZ ¼ IZo þ Ix þ Iyi ½sinðh0 þ h Þ 2 þ Izi ½cosðh0 þ h Þ 2 IZo þ Ix þ Iyi ðsin h0 þ h

cos h0 Þ2 þ Izi ðcos h0 h sin h0 Þ2 IZo þ Ix þ Iyi ðsin2 h0 þ 2h sin h0 cos h0

þ 0Þ þ Izi ðcos2 h0 2h sin h0 cos h0 þ 0Þ ¼ IZo þ Ix þ Iyi sin2 h0 þ Izi cos2 h0
þ 2h ðIx þ Iyi Izi Þ sin h0 cos h0 ¼ I0 þ 2h ðIx þ Iyi Izi Þ sin h0 cos h0
½I0 þ 2h ðIx þ Iyi Izi Þ sin h0 cos h0 wþ Iz nðcos h cos h0 Þ I0 wþ 2h w
ðIx þ Iyi Izi Þ sin h0 cos h0 þ Iz nðcos h0 h sin h0 cos h0 Þ ¼ 0

ðIx þ Ixi Þhþ Izi Ix Iyi w2 sin h cos h þ Iz wn sin h ðIx þ Ixi Þh þ Iz wn

ðsin h0 þ h cos h0 Þ þ Izi Ix Iyi w2 ½sin h0 cos h0 þ h ðcos2 h0 sin2 h0 Þ

¼ ðIx þ Ixi Þh þ Iz wn sin h0 þ fIz wnh cos h0 þ Izi Ix Iyi w2 sin h0 cos h0

þ Izi Ix Iyi w2 h ðcos2 h0 sin2 h0 Þ ¼ ðIx þ Ixi Þh þ Iz wn sin h0
þ fIz wnh cos h0 ðIx þ Iyi Izi Þ½sin h0 cos h0 þ h ðcos2 h0 sin2 h0 Þ w2 g ¼ 0
The second and the third of the three equations written above are nonlinear.
Thomson [10] has solved them by using the perturbation technique, which is
described by means of the following example.
Let us consider the following nonlinear differential equation of the first order
xþ ax þ bx2 ¼ 0
where x is an unknown function of the independent variable t, and the coefficient

b is a small quantity in comparison with the other coefficients 1 and a.
Let us consider another equation similar to the first
xþ ax þ lbx2 ¼ 0
which contains an additional factor l which may be any positive number. If the
solution x of the second equation were known, then the solution of the first would
also be known by setting l = 1. For this purpose, we seek a solution of the second
equation in the following form
x ¼ x0 þ lx1 þ l2 x2 þ
where x0, x1, …, xn are unknown functions of the independent variable t.

By substituting x = x0 + lx1 + l2x2 + into x` + ax + lbx2 = 0, there results
2
x0 þ lx1 þ l2 x2 þ þ a x0 þ lx1 þ l2 x2 þ þ lb x0 þ lx1 þ l2 x2 þ
¼0
The preceding equation may also be written as follows
ðx0 þ ax0 Þ þ lðx1 þ ax1 þ bx0 2 Þ þ l2 ðx2 þ ax2 þ bx0 x1 Þ þ l3 ðx3 þ Þ ¼ 0
In case of l = 0, the function x0 is the solution of the linear equation

x` + ax = 0. The function x0 is called the generating solution, which must
satisfy the initial condition of the problem.
In case of l 6¼ 0, the preceding equation
ðx0 þ ax0 Þ þ lðx1 þ ax1 þ bx0 2 Þ þ l2 ðx2 þ ax2 þ bx0 x1 Þ þ l3 ðx3 þ Þ ¼ 0
can be satisfied only if the coefficients of l, l2, l3, … are equal to zero, that is, if
x1 þ ax1 þ bx20 ¼ 0

x2 þ ax2 þ bx0 x1 ¼ 0
x3 þ ¼ 0
These equations are to be solved for the functions x1, x2, …, and so on.
This technique has been applied by Thomson [10] to the two nonlinear differ-
ential equations found previously and rewritten below for convenience
I0 w ðIz n sin h0 Þh þ 2h wðIx þ Iyi Izi Þ sin h0 cos h0 ¼ 0

ðIx þ Ixi Þh þ ðIz n sin h0 Þwþ fIz wnh cos h0 ðIx þ Iyi Izi Þ½sin h0 cos h0
þ h ðcos2 h0 sin2 h0 Þ w2 g ¼ 0
where the solution has been computed only to the first-order correction. Since the
symbol h0 has been used in the two equations written above, then the symbols h00
and w`00 are used below for the solutions of the two linear differential equations
(corresponding to the function x0).
These linear equations are written below.
I0 w00 ðIz n sin h0 Þh00 ¼ 0

ðIx þ Ixi Þh00 þ ðIz n sin h0 Þw00 ¼ 0
Solving for w`00, there results

" #
ðIz n sin h0 Þ2
h00 ¼ h00 ¼ 0
I0 ðIx þ Ixi Þ
After setting for convenience
ðIz n sin h0 Þ2
x2 ¼
I0 ðIx þ Ixi Þ
the generating solution which satisfies the initial conditions can be written as
follows

h0
h00 ¼ sinðxtÞ
x

Iz n sin h0 Iz n sin h0 h0
w00 ¼ h00 ¼ sinðxtÞ
I0 I0 x
Now we consider the first-order (x1) correction to the generating solution (x0).
This solution makes it possible to compute the following nonlinear terms
Iz nh20 sin h0 2
h00 w00 ¼ sin xt
I0 x 2
2
h0 Iz n sin h0 2 2
w200 ¼ sin xt
x I0
These terms are to be substituted into the two nonlinear differential equations
I0 w ðIz n sin h0 Þh þ 2h wðIx þ Iyi Izi Þ sin h0 cos h0 ¼ 0

ðIx þ Ixi Þh þ ðIz n sin h0 Þwþ fIz wnh cos h0 ðIx þ Iyi Izi Þ½sin h0 cos h0
þ h ðcos2 h0 sin2 h0 Þ w g ¼ 0
2
By so doing, Thomson [10] obtains

I0 w1 ðIz n sin h0 Þh1 ¼ 2½ðIx þ Iyi Izi ÞIz nh0 2 sin2 h0 cos h0 Þ= I0 x2 sin2 xt

ðIx þ Ixi Þh1 þ ðIz n sin h0 Þw1 ¼ ðIz n cos h0 Þ½ðIz nh0 2 sin h0 Þ= I0 x2 sin2 xt
þ ðIx þ Iyi Izi Þ sin h0 cos h0 ½ðIz nh0 sin h0 Þ=ðI0 xÞ 2 sin2 xt
By solving the first of the two equations written above for h1 and differentiating
twice h1 with respect to time, Thomson [10] obtains the following expression
h1 ¼ ½I0 =ðIz n sin h0 Þ w1 þ f½4h0 2 ðIx þ Iyi Izi Þ cos h0 sin h0 =½I0 gðcos2 xt
sin2 xtÞ
By substituting this expression of h``1 into the second of the two equations
written above, Thomson [10] obtains the following differential equation
½I0 ðIx þ Ixi Þ=ðIz n sin h0 Þ w1 þ ðIz n sin h0 Þw1 ¼ ðh0 Iz n=xÞ2 f½ðIx þ Iyi Izi Þ sin3 h0
cos h0 =½I0 2 ½sin h0 cos h0 =½I0 g1=2½1 cosð2xtÞ ½ðIx þ Ixi Þ=I0 ½4h0 2
ðIx þ Iyi Izi Þ sin h0 cos h0 cosð2xtÞ
As to this differential equation, the solution of the homogeneous equation for w`1
is also harmonic of angular frequency x = (Izn sin h0)/[I0(Ix + Ixi)]1/2.
The particular solution has harmonic terms of angular frequency 2xt. In addi-
tion, it has a constant term equal to the constant term on the right-hand side of the
equation divided by the coefficient of w`1 on the left-hand side.
This constant term results in a steady drift which rotates the outer gimbal
according to the equation ws = w`s t. The constant term of the solution is [1]:
ws ¼ ½h0 2 Iz n=ð2x2 I0 2 sin h0 Þ ½ðIx þ Iyi Izi Þ sin3 h0 cos h0 I0 sin h0 cos h0
Remembering that I0 = IZo + (Ix + Iyi) sin2h0 + Izi cos2h0, the preceding
expression can also be written as follows
h20 Iz nðIZo þ Izi Þ cos h0

ws ¼
2x2 I02
The minus sign in front of the term on the right-hand side of the preceding
expression indicates that the outer gimbal oscillates and drifts in a negative
direction. This phenomenon is known as gimbal walk. In other words, under the
conditions specified below, the motion of the outer gimbal of a free gyroscope is not
purely oscillatory but rather shows some drift. This drift is caused by the moments
of inertia of the gimbals. The preceding expression shows that there is no gimbal
walk for either h0 = p/2 or IZo + Izi = 0, because in both cases w`s = 0.
8.7 The Gyrocompass
A gyrocompass is a north-seeking gyroscope carried on board a vehicle, for the

purpose of indicating the direction of true north at any time, let the latitude of the
carrying vehicle be what it may. To this end, the rotor of the gyroscope is con-
tinuously driven, and its spin axis is parallel to the axis of rotation of the Earth.
Following Bowditch [12], in order for a gyroscope to act as a gyrocompass, the
rotor is mounted in a sphere, which in turn is supported in a vertical ring. This
assembly is mounted on a base, as shown in the following figure, which is due to
the courtesy of Sperry Gyroscope Company [13].
A gyrocompass has two advantages over a magnetic compass:

• it finds true north, which is determined by the rotation of the Earth, instead of
magnetic north; and
• it is insensitive to ferromagnetic materials, which change the magnetic field
around them.
8.7 The Gyrocompass 1081
By true north, or geodetic north, we mean the direction along the surface of the
Earth towards the geographic north pole, that is, the direction of the local meridian.
A gyrocompass is based on two properties of a gyroscope, which have been shown
at length in the preceding sections. They are gyroscopic inertia (which is the ten-
dency of any rotating body to preserve its plane of rotation, in accordance with
Newton’s first law of motion) and precession (a gyroscope rotates not only around
its own axis but also around the vertical axis).
These properties of gyroscopes are used in a gyrocompass together with two
natural phenomena, which are the rotation of the Earth about its axis and the force
of gravity, as will be shown below. As a result of these facts, a gyrocompass on
board a vehicle aligns itself with the local meridian, and therefore indicates con-
stantly the direction of true north, independently of the motions of yaw, pitch, and
roll of the carrying vehicle.
If a spinning gyroscope were placed on land at the terrestrial equator, with the
axle of the gyroscope parallel to the polar axis of the Earth, then the axle would
remain aligned with the direction of the local meridian, because of the absence of
any force tending to deflect the axle from this direction.
On the other hand, when a gyroscope is carried on board a vehicle, then its axle
is desired to remain (in order for the instrument to be of any use) aligned with the
local meridian at any latitude, independently of the motion of the vehicle. To this
end, a gyroscope must be made in such a way as to seek and keep the direction of
the meridian against the friction of its supports and other perturbing forces due to
the motion of the vehicle.
In order to seek and maintain true north, a gyroscope must:
• be made to stay in the plane of the local meridian;
• be made to remain horizontal; and
• stay in this position after reaching it independently of the motion of the vehicle.
To seek the plane of the local meridian, a weight is added to the bottom of the
vertical ring, so as to cause it to swing on its vertical axis, and thus seek to align
itself horizontally. Since the gyroscope tends to oscillate, a second weight is added
to the side of the sphere in which the rotor is contained, in order to dampen the
oscillations until the gyroscope stays on the meridian. Due to these two weights, the
only possible point of equilibrium for the gyroscope is on the meridian with its spin
axis horizontal.
In order to seek north, a gyroscope is provided with a system of reservoirs
partially filled with mercury. This system of reservoirs, known as mercury ballis-
tics, is used to apply a force against the spin axis of the rotor. The mercury
ballistics, which are usually four, are placed so that their centres of gravity coincide
exactly with the centre of gravity of the gyroscope. The flow of mercury causes the
tip of the spin axis to trace an ellipse instead of a circumference in the motion of
precession, as shown on the left-hand side of the following figure, which is also due
to the courtesy of Sperry Gyroscope Company [13]. Each of these ellipses takes
about 84 min to complete. This value results from the well-known formula
12
‘
T ¼ 2p
g
which expresses the period T of a pendulum whose length ‘ is equal to the radius of
a spherical Earth. To dampen the motion of precession, the force is applied not in
the vertical plane, but slightly to the east of the vertical plane. This causes the tip of
the spin axis to trace a spiral instead of an ellipse, and finally settle on the meridian
pointing north, as shown on the right-hand side of the following figure.
With reference to the following figure, the working principle of a gyrocompass

can be mathematically described as follows.
The double-gimbal gyroscope illustrated above has a pendulous mass m, whose

weight is w = mg, placed at a distance ‘ from the centre of mass O of the rotor
along the y-axis, in order for the gyroscope to have a moment
MOx ¼ w‘ cos h
about the x-axis when the z-axis forms an angle a above the horizontal plane.
With reference to the following figure, let xE be the angular velocity vector due
to the rotation of the Earth from west to east about its polar axis. This vector points
north. For a given point placed at a geocentric latitude u* on the surface of a
spherical Earth, the radial component and the tangential component of the vector
xE are, respectively, xE sin u* and xE cos u*.
Let the Z-axis of a gyrocompass be oriented along the local vertical direction. In
order for the z-axis to remain in the plane of the local meridian (i.e. in order for the
z-axis to point north), the outer gimbal must have a steady motion of precession at a
rate w` = xE sin u* and in addition must have an angular velocity equal to xE cos u*
about the y-axis perpendicular to the plane of the outer gimbal. The gyroscope is
assumed to be constrained to move in this manner.
We want to determine the moment MOx necessary for this motion.
Let a0 be a small angle which the z-axis forms with respect to the horizontal
plane when the given point is at an angle u* of geocentric latitude. In this case, the
Eulerian angle h between the axes z and Z is h = p/2 − a0, and the components
along the axes x, y, and z of the angular velocity vector xE are
xEx ¼ 0
xEy ¼ xE sinðu a0 Þ
xEz ¼ xE cosðu a0 Þ
When the spin rate /` of the rotor is very large, the moment-of-momentum
vector hO can be assumed to have only the component hOz along the z-axis.
Since MOz = 0, then /` = constant. Therefore, the required moment MOx to be

applied about the x-axis by means of the weight w is
MOx ¼ ðIz /ÞxEy
Since
MOx ¼ w‘ cos h ¼ w‘ sin a0 and xEy ¼ xE sinðu a0 Þ
then
w‘ sin a0 ¼ Iz /xE sinðu a0 Þ ¼ Iz /xE ðsin u cos a0 cos u sin a0 Þ
By dividing all terms of the preceding equation by sin a0 and solving for tan a0,
there results
Iz /xE sin u
tan a0 ¼
w‘ þ Iz /xE cos u
which expresses the tangent of the required angle a0 of inclination of the spin axis
of the gyrocompass above the horizontal plane. This angle also depends on the
geocentric latitude u* of the given point on the surface of the Earth.
The moment required for the angular velocity xE cos u* about the y-axis is due
to the reaction of the bearings placed along the Z-axis of the outer gimbal.
Let us consider a gyrocompass whose spin axis z is deviated by a perturbation
from the plane of the local meridian, as shown in the following figure.
The angle of deviation resulting from this perturbation has two components. One
(w) of them is perpendicular to the meridian plane, and the other component (a) is
contained in the meridian plane. Since the two components have the same fre-
quency, then the spin axis z of the rotor has a motion of precession and its tip
describes an ellipse, as has been shown above.
When both of these angular components have small values, then the angular
velocities about the axes x, y, and z are respectively
xx ¼ a xE w cos u
xy ¼ wþ xE sin u xE a cos u
xz ¼ ðwþ xE sin u Þa þ xE cos u
As has been done above, so here we assume the spin rate /` of the rotor to be
very large, so that the moment-of-momentum vector hO has only the component hOz
along the z-axis, because the other two components hOx and hOy of hO are negligible
in comparison with hOz.
In other words, for MOz = 0, there results hOz = Iz/` = constant, and
hOx = hOy = 0.
The equations expressing the moments about the axes x, y, and z are respectively
MOx ¼ hOz xy ¼ Iz /ðwþ xE sin u xE a cos u Þ ¼ w‘a

MOy ¼ hOz xx ¼ Iz /ðaþ xE w cos u Þ ¼ 0
MOz ¼ 0
where w‘a = mg‘a is the moment due to the weight acting on the mass m placed
along the negative direction of the y-axis at a distance ‘ from the centre of mass O
of the rotor. The first equation can also be written as follows
ðIz /Þw ðIz /xE cos u þ w‘Þa ¼ Iz /xE sin u
The second equation yields
aþ xE w cos u ¼ 0
This equation, solved for w, yields
a
w¼
xE cos u
which in turn, differentiated with respect to time, yields
a
w ¼
xE cos u
This expression of w` substituted into the first equation rewritten below
ðIz /Þw ðIz /xE cos u þ w‘Þa ¼ Iz /xE sin u

leads to the following differential equation
ðIz /xE cos u þ w‘ÞxE cos u

aþ a ¼ xE 2 sin u cos u
Iz /
The general solution of this equation is
Iz /xE sin u
a ¼ c1 sin pt þ c2 cos pt þ
Iz /xE cos u þ w‘
where
1 1
ðIz /xE cos u þ w‘ÞxE cos u 2 w‘xE cos u 2
p¼
Iz / Iz /
Remembering the preceding equation w = −a`/(xE cos u*), the solution for w is
p
w¼ ðc1 cos pt c2 sin ptÞ
xE cos u
In other words, the spin axis z of the rotor oscillates horizontally about the plane
of the local meridian by the angle w and also oscillates vertically about the plane
which forms the stationary angle a0 with respect to the equator.
The angular frequency p [(w‘xE cos u*)/(Iz/`)]1/2 of this oscillation is a
function of the geocentric latitude u*. The value of this angular frequency is very
small, due to the high value of Iz /`. For values of u* approaching p/2, the value of
p approaches zero, and therefore the reliability of a gyrocompass decreases at high
values of latitude.
The results found above are based on the assumption hOx = hOy = 0. In such
conditions, a gyrocompass has only the low angular frequency p of oscillation.
In case of hOx and hOy being not negligible in comparison with hOz, it is nec-
essary to modify the two preceding equations rewritten below
MOx ¼ Iz /ðwþ xE sin u xE a cos u Þ ¼ w‘a

MOy ¼ Iz /ðaþ xE w cos u Þ ¼ 0
by adding the terms −Ix a`` and Iy w`` to, respectively, the first and the second of
them. The third equation MOz = 0 does not change. By so doing, there results
MOx ¼ Iz /ðwþ xE sin u xE a cos u Þ Ix a ¼ w‘a

MOy ¼ Iz /ðaþ xE w cos u Þ þ Iy w ¼ 0
After setting for convenience
a ¼ Iz /xE cos u
b ¼ w‘ þ Iz /xE cos u
the preceding differential equations can be rewritten as follows

Iz / b Iz /xE sin u
a wþ a¼
Ix Ix Ix

Iz / a
wþ aþ w¼0
Iy Iy
The angular oscillations are assumed to be harmonic, that is, a = Aa exp(ipt) and
w = Aw exp(ipt), where i = (−1)1/2. In this case, the values p of the natural angular
frequencies of these oscillations result from solving the equation
detðAÞ ¼ 0
where A is the following 2 2 matrix

p2 þ b=I ðIz /=Ix Þip
A¼ x
Iz /=Iy ip p2 þ a=Iy
By expanding the determinant of the matrix A and equalling det(A) to zero, there
results
" 2 #
aI x þ bI y þ Iz / ab
p4 p2 þ ¼0
Ix Iy Ix Iy
Since aIx and bIy are, both of them, negligible in comparison with (Iz /`)2, then
the preceding equation reduces to
" #
ðIz /Þ2 2 ab
p 4
p þ 0
Ix Iy Ix Iy
The two solutions p21 and p22 of the preceding equation are
" #12 8 " #12 9
ðIz /Þ2 1 ðIz /Þ4 4ab ðIz /Þ2 < 4abIx Iy =
p21;2 ¼ ¼ 1 1
2Ix Iy 2 Ix2 Iy2 Ix Iy 2Ix Iy : ðIz /Þ4 ;
Since the numerator of the fraction 4abIxIy/(Iz/`)4, which is on the right-hand

side of the preceding expression, is much smaller than the denominator, then the
function [1–4abIxIy/(Iz/`)4]1/2 may be expanded in a Maclaurin series, and this

expansion may be truncated after the first derivative. By so doing, the preceding
expression becomes
8 " #12 9 ( " #)
ðIz /Þ2 < 4abIx Iy = ðIz /Þ2 2abIx Iy
p1;2 2 ¼ 1 1 1 1
2Ix Iy : ðIz /Þ4 ; 2Ix Iy ðIz /Þ4
Let p21 and p22 be the squares of the two (low and high) angular frequencies
corresponding, respectively, to the minus sign and to the plus sign placed in front of
[1–2abIxIy/(Iz/`)4].
Remembering that a = Iz/`xE cos u* and b = w‘ + Iz/`xE cos u*, there results
ab ½w‘ þ ðIz /ÞxE cos u xE cos u

p1 2 ¼ 2
¼
ðIz /Þ Iz /
The low angular frequency p1 turns out to have the same value as that computed
previously in case of hOx = hOy = 0.
As to the high angular frequency p2, there results
" #
ðIz /Þ2 abIx Iy ðIz /Þ2
p2 2
¼ 1
Ix Iy ðIz /Þ4 Ix Iy
In other words, when the assumption hOx = hOy = 0 is valid, then a gyrocompass
oscillates at only one angular frequency, whose value is p1. Otherwise, there are two
values, of which one (p1) is low and the other (p2) is high, of the angular frequency
of a gyrocompass. However, the amplitude of the oscillation at the high angular
frequency is very small. Therefore, the oscillation at the low angular frequency is
the only one which can be perceived in practice.
As has been shown above, the oscillation of a gyrocompass at the low angular
frequency is damped by a moment generated by the weight force w = mg acting on
a point-mass m placed along the negative direction of the y-axis at a distance ‘ from
the centre of mass O of the rotor. This point-mass is also placed at a distance e to
the east of the y-axis, so that its co-ordinates in the body-fixed system of reference
are x = −e, y = −‘, and z = 0. In order to take account of the damping moment
MOy = −wea, the preceding equations which express the moments about the axes
x, y, and z

MOy ¼ Iz /ðaþ xE w cos u Þ ¼ 0
MOz ¼ 0
are modified as follows

MOy ¼ Iz /ðaþ xE w cos u Þ ¼ wea
MOz ¼ 0
The second of the three equations written above, solved for a`, yields
wea
a ¼ xE w cos u
Iz /
The first of the three equations, solved for a, yields
Iz /ðwþ xE sin u Þ
a¼
Differentiating this function with respect to time, there results
Iz /w
a ¼
By substituting these expressions of a and a` into the equation

wea
a ¼ xE w cos u
Iz /
there results
we ðIz /xE cos u þ w‘ÞxE cos u wexE sin u

wþ wþ wþ ¼0
Iz / Iz / Iz /
The preceding differential equation shows that the oscillation w of a gyrocom-

pass is damped in amplitude by the offset e of the weight w = mg. This offset moves
the equilibrium point of the instrument to the east. The angle w0 through which the
equilibrium point is moved results from setting w`` = w` = 0 in the preceding
equation and solving for w. This yields
we tan u
w0 ¼
A gyrocompass is subject to errors. They form combined together the so-called

gyro error, which is expressed in degrees east or west. According to Bowditch [3],
the gyro error is constant in one direction; that is, an error of (for example) one
degree east applies to all bearings all around the compass.
The errors to which a gyrocompass is subject are speed error, latitude error,
ballistic deflection error, ballistic damping error, quadrantal error, and gimballing
error. Further error may be caused by a malfunction of the instrument or its
incorrect alignment with the central line of the carrying vehicle.
Speed error is due to the fact that a gyrocompass only moves directly east or
west when it is either stationary on the rotating Earth or placed on a vehicle moving
exactly east or west on the surface of the Earth. Therefore, any motion of the
vehicle to the north or south causes a gyrocompass to trace a path which depends on
the speed of this motion and also on the amount of northerly or southerly heading,
as will be shown below. This fact causes a gyrocompass to settle a bit off true north.
The resulting error is westerly if the course of the carrying vehicle is northerly, and
easterly if this course is southerly.
Let us consider first the case of a vehicle which carries a gyrocompass and also
moves northerly along the local meridian at a velocity v. This motion of the vehicle
generates an angular velocity vector v/rE pointing west, where rE is the radius of a
spherical Earth. This vector, combined with the component xE cos u* (which lies
on the plane of the horizon) of the angular velocity vector xE of the Earth, causes
the resultant angular velocity in the plane of the horizon to deviate to the west by an
angle c. An approximate value in radians of this angle is
v
rE
c
xE cos u
Because of this deviation, the axle of the gyrocompass points in the direction of
the resultant angular velocity vector, the angle c being the heading error.
With reference to the following figure, let us consider now the case of another
vehicle moving at a velocity v which forms an angle h with the local meridian.
In this case, the projections of the velocity vector v of the vehicle onto,
respectively, the local meridian and the local parallel are v cos h and v sin h.
Therefore, the preceding equation becomes
v cos h
rE v cos h
c ¼
v sin h rE xE cos u þ v sin h
xE cos u þ
rE
In practice, the effect of the term v sin h can be neglected, because this term is
small in comparison with rExE cos u*. The second of the two cases reduces to the
first when h = 0°. As has been shown above, the magnitude of this error depends on
the speed (v), course (h), and latitude (u*) of the vehicle.
This error can be corrected internally by means of a cosine cam. This mechanism
is substantially a roller which rides in a cosine groove which is cut into the lower
side of the azimuth gear. The cosine cam is designed to move the corrector
mechanism as necessary to correct the compass readings for all changes in course of
the vehicle. A description and a drawing of this mechanism can be found in Ref.
[14].
Tangent latitude error affects only gyrocompasses with mercury ballistics and is
easterly in north latitudes and westerly in south latitudes. This error, too, is cor-
rected internally, by offsetting the lubber line, or by means of a small movable
weight attached to the casing. By the way, in the window of the compass is a line
set parallel with the fore-and-aft line of the ship. This is called the lubber line. The
compass is read by noting the markings directly behind this line [15].
Ballistic deflection error occurs when there is a marked change in the north–
south component of the speed of the vehicle. East–west accelerations have no
effect. A change of course or speed of the vehicle also results in speed error in the
opposite direction, so that the two errors tend to cancel each other if the gyro-
compass is properly designed. This aspect of design involves a slight offset of the
ballistics according to the operating latitude, upon which the correction depends. As
latitude changes, the error becomes apparent, but can be reduced by adjusting the
offset.
Ballistic damping error is a temporary oscillation due to changes in course or
speed of the vehicle. During one of such changes, the mercury in the ballistics is
subject to centrifugal and acceleration/deceleration forces. These forces generate
torques about the spin axis end therefore errors in the gyrocompass reading. Slow
changes do not introduce significant errors, but rapid changes do. In order to
counteract this cause of error, the position of the ballistics is changed so that the true
vertical axis is centred, and therefore is not subject to error. This is done only when
certain rates of turn or acceleration are exceeded.
Quadrantal error is due to two causes. The first is the not perfect centring of the
centre of gravity in the outer frame (the so-called follow-up or phantom system).
Due to this cause, the gyroscope tends to swing along its heavy axis as the vehicle
rolls. For the purpose of counteracting this cause of error, some weight is added, so
that the mass is the same in all directions from the centre. There being no long axis
of weight, there is also no tendency of the gyroscope to swing in one particular
direction. The second cause of quadrantal error is the rolling of the vehicle. As the
vehicle rolls, the apparent vertical axis is displaced, first to one side and then to the
other. The vertical axis of the gyroscope tends to align itself with the apparent
vertical axis. In case of northerly or southerly courses, and also in case of easterly or
westerly courses, the precession of the gyrocompass is equal on both sides, so that
the resulting error is zero. In case of intermediate courses, the N–S and E–W
precessions are additive, so that a persistent error is introduced, which changes
direction in different quadrants. This error is corrected by means of a second
gyroscope, called the floating ballistics, which stabilises the mercury ballistics as
the vehicle rolls. Another method consists in using two gyroscopes for the directive
element. The motion of precession of these gyroscopes occurs in opposite direc-
tions, so as to eliminate this error.
Gimballing error is caused by readings taken from the compass card when it is
tilted from the horizontal plane. This applies to the compass itself and to all
repeaters. To reduce this error, the outer ring of the gimbal of each repeater should
be installed with the fore-and-aft line of the vehicle. Of course, the lubber line must
also be exactly centred.
8.8 The Rate Gyroscope
Rate gyroscopes, also called rate-of-turn gyroscopes, are instruments used for
guidance and control of vehicles, in order to detect and measure rates of change of
angles. They consist essentially of a spinning rotor supported by an inner gimbal,
which in turn is restrained by a spring mechanism. This mechanism permits a
limited rotation of the inner gimbal about the outer gimbal, which is fixed to the
vehicle. The axis (the Z-axis shown in the following figure) about which the vehicle
turns is called the input axis.
8.8 The Rate Gyroscope 1093
The axis of rotation of the inner gimbal is called the output axis. Let Iz and n be,
respectively, the moment of inertia and the angular velocity of the rotor about its
spin axis (z). The magnitude hO of the moment-of-momentum vector hO is
hO ¼ hOz ¼ Iz n
Let us suppose that the vehicle makes a steady turn about the input axis Z, at a
rate w`. As a result of this turn, the moment-of-momentum vector changes at a rate
Iznw`. This change requires a moment of the same magnitude about the output axis.
As the inner frame rotates by a small angle h about the output axis, this moment is
provided by a torsional spring of stiffness k.
By equalling the two moments, there results
Iz nw ¼ kh
which, solved for h, yields
Iz n
h¼ w
k
The preceding equation shows that the angle h, through which the inner frame
rotates about the output axis, is proportional to the rate-of-turn w` of the outer frame
about the input axis. Since the outer frame is fixed to the vehicle, then this angle is
proportional to the rate of turn of the vehicle. In the simple scheme illustrated
above, the turn made by the vehicle has been supposed to occur in the counter-
clockwise direction. However, the vehicle may also turn in the opposite (clockwise)
direction. Therefore, the spring mechanism must be so designed as to restrain the
inner gimbal from rotating in either direction.
The angular displacement h is detected by an E-pickoff device, shown in the

preceding figure (redrawn from [1]), which converts this angle to a voltage pro-
portional to the angle itself. The output of the E-pickoff device is the rate signal.
An alternating current, whose frequency is generally 400 Hz, is applied to the
middle leg of the E-shaped device. The windings on the other two legs are in
opposition. Therefore, when the armature, which is attached to the output axis, is
centred with respect to the middle leg, then no voltage is detected across the
opposed outer coils, which are connected in series. By contrast, when the armature
is displaced from its central position by an angle h, then the path of the magnetic
flux is unbalanced, and a voltage is sensed by the instrument across the outer coils.
If a damping apparatus were absent, the output axis would first overshoot the
angular displacement h, and then oscillate about it. In order to prevent these
undesired oscillations about the output axis, the inner gimbal has generally a
damping system. The following figure, due to the courtesy of NASA [16], illus-
trates a rate gyroscope, whose inner gimbal is restrained from rotating by either an
elastic mechanism or a torquer, and which is also equipped with a damping system.
8.8 The Rate Gyroscope 1095
Another type of rate gyroscope is the so-called floated gyro unit, which uses a
torsion bar instead of a spring system, as shown in the following figure (redrawn
from Ref. [17]). The torsion bar is mounted along the output axis of the gyroscope
and produces restraining torques in either rotational direction. The inner gimbal
sphere turns about the output axis in response to a rotation about the input axis. The
case contains a fluid which surrounds the inner gimbal sphere and provides flota-
tion. The same fluid also damps the oscillations due to sudden changes in the
angular rate input and provides protection from shocks.
The behaviour of a rate gyroscope is governed by the following differential

equation
Ix hþ chþ kh ¼ Iz nw
where Ix is the moment of inertia of the system comprising the rotor and the inner
gimbal about the output axis, and c is the coefficient of viscous damping.
The characteristics of the rate gyroscope can determined by considering the

behaviour of the instrument in the transient period (without the forcing function
Iznw`). This behaviour is governed by the corresponding homogeneous equation
hþ 2fxhþ x2 h ¼ 0
where x = (k/Ix)1/2 is the natural angular frequency of the instrument in the absence
of damping, f = c/ccr is the damping ratio, and ccr = 2(kIx)1/2 is the critical
damping.
8.9 The Rate Integrating Gyroscope
In a rate integrating gyroscope, the spring system is removed and replaced with a
high-precision ball bearing. A scheme of this type of gyroscope is shown in the
following figure, due to the courtesy of NASA [16], where IA, OA, and SA stand
for, respectively, input axis, output axis, and spin axis.
The only restraint existing in a rate integrating gyroscope is due to the damping
fluid. When a torque arises about the gimbal axis as a result of an angular velocity
about the input axis, then the gimbal begins to rotate about its own (or output) axis.
As shown above, there being no spring system, the only thing which opposes the
rotation of the gimbal is the presence of the damping fluid. The resistance opposed
by this fluid to the gimbal rotation depends on its viscosity and also on the rota-
tional velocity of the gimbal.
In mathematical terms, the behaviour of a rate integrating gyroscope is governed
by the same differential equation as that of the preceding section
Ix hþ chþ kh ¼ Iz nw
where the terms Ixh`` and kh are both of them equal to zero, because h` is constant
in time and there is no spring system, h` and w` are the angular velocities about,
respectively, the output axis (x) and the input axis (Z), Ix is the moment of inertia of
8.9 The Rate Integrating Gyroscope 1097
the system comprising the rotor and the inner gimbal about the output axis, Iz and
n are, respectively, the moment of inertia and the angular velocity of the rotor about
its spin axis (z), and c is the coefficient of viscosity of the damping fluid. In these
conditions, the preceding equation becomes
ch ¼ Iz nw
which, solved for h` and integrated, yields
Zt
Iz n Iz n
h¼ wds ¼ w
c c
0
In other words, the angle h expressing the rotation about the output axis is
proportional to the integral of the angular velocity w` about the input axis, that is, h
is proportional to the angle w of rotation about the input axis. Therefore, by means
of a pickoff device, it is possible to measure the angle w through which the
gyroscope has rotated about its input axis.
A rate integrating gyroscope also contains an electrical torquer which causes the
gimbal to have a motion of precession until the pickoff signal goes back to zero.
This type of gyroscope operates in a closed loop system, where the pickoff output
signal drives a servo amplifier, which supplies current to the electrical torquer. By
so doing, what is actually read is not the pickoff angle, but rather the amount of
current supplied to the electrical torquer. The current coming from the servo
amplifier is proportional to the torque which acts on the gimbal [18].
8.10 High-Precision Gyroscopes
Some space missions require high-precision measurements of angular velocities.

One of such missions is the Gravity Probe B (GP-B) experiment [19], whose
purpose is to test two phenomena, namely the geodetic effect and the
frame-dragging effect, predicted by Einstein’s general theory of relativity. To this
end, it is necessary to perform measurements with great precision, because the
geodetic effect causes the orientation of the spin axis of a gyroscope carried by a
satellite circling the Earth along a polar orbit to change by 6.6 arcseconds (or
0.0018 degrees) in a year with respect to a distant star taken as a point of reference.
Likewise, the frame-dragging effect causes the orientation of the spin axis of the
same gyroscope to change in the plane orthogonal to the orbital plane of the Earth
by 0.041 arcseconds (or 0.000011 degrees) in a year.
According to the laws of classical (or Newtonian) physics, a perfect gyroscope,
which is not subject to external forces, does not drift, that is, the orientation of its
spin axis continues to stay aligned with a direction of reference, which is the
direction of a distant star. By contrast, in accordance with Einstein’s general theory
of relativity, the geodetic and frame-dragging effects cause a slight change (by the
amounts indicated above) of this orientation.
The principal components of the science instrument of the GP-B experiment are
four gyroscopes, the optical telescope, and the mounting block. All of these
components are made of fused quartz, because of the high stability (very little and
uniform expansion and contraction) of this material over a wide range of temper-
atures. The guide star for the GB-P experiment is IM Pegasi, which is a variable
binary star system approximately 329 light years away in the constellation of
Pegasus. This star, visible to the naked eye, was chosen because its microwave
radio emissions are observable with a large radio telescope network on the ground,
so that its precise position can be related by interferometry to distant quasars.
The rotors used for the GB-P experiment are perhaps, of all the objects made so
far, those which best approximate the shape of a sphere. Such rotors are shown in
the following figure, which is due to the courtesy of the Stanford University [19].
To give an idea of the degree of perfection reached by them, suffice it to say that,
if the rotors, whose size is about the same as that of a ping pong ball, were enlarged
to the size of the Earth, then the tallest mountain or the deepest valley would be
only 2.4 m in height. The gyroscopes are housed in a large Dewar (thermos like)
reservoir containing 2441 litres of superfluid helium, maintaining a temperature less
than 2 K (−271 °C). Such a temperature, which is near to absolute zero, is nec-
essary to reduce molecular interference to a minimum and also to enable the lead
and niobium components of the gyroscope mechanism to become superconductive.
Superconductivity is necessary to protect the gyroscopes from the magnetic field
of the Earth, because a faint magnetic signal from the gyroscopes is used to detect
the change in angle of their spin axes. The intrusion of Earth’s magnetic field would
swamp that signal.
This extremely low temperature also makes it possible to create an ultra-low
pressure vacuum in the gyroscope chamber. After pumping out most of the gas, the
molecules of gas that remain are very cold and thus hardly moving, which means
that they exert almost zero pressure. In this high-vacuum environment, the four
8.10 High-Precision Gyroscopes 1099
spherical gyroscopes spin at their operating speed of about 72 Hz (or 4300 rpm).
The rotors are of spherical shape, because a superconducting sphere, when spun,
produces a weak magnetic field which is precisely aligned with its axis of rotation.
The gyroscopes are therefore coated with a metallic layer of niobium of near-perfect
uniformity. At the cryogenic temperature in the core of GP-B, niobium becomes a
superconductor and therefore produces a magnetic field when the spheres are spun.
By monitoring the magnetic field, it is possible to monitor the spin of the gyro-
scopes. The GP-B instrument is designed to measure changes in the orientation of
the spin axes of the gyroscopes to better than 0.5 arcseconds (or 1.4 10−7
degrees) over a period of one year.
As has been shown in Sect. 3.24, the analysis of the data coming from all of the
four gyroscopes indicated a geodetic drift rate of −6601.8 ± 18.3 mas/year and a
frame-dragging drift rate of −37.2 ± 7.2 mas/year. These values are to be com-
pared with the general relativity predictions of −6606.1 and −39.2 mas/year,
respectively, where 1 mas = 4.848 10−9 rad [20].
8.11 Optical Gyroscopes
The gyroscopes described hitherto are, all of them, rotary gyroscopes, which are
based on the principle of conservation of the moment–of-momentum vector (in
magnitude and direction) in the absence of resulting torques.
There are gyroscopes based on other physical principles, which can be classified
into the two categories of optical gyroscopes and vibrating structure gyroscopes.
Those belonging to the first category are briefly described in the present section.
Optical gyroscopes are instruments having no moving parts, which are used for the
purpose of measuring angular velocities. Two common types of optical gyroscopes
are fibre optic gyroscopes and ring laser gyroscopes, which are shown in the present
section. They are instruments which sense changes in orientation and therefore
perform the functions of rotary gyroscopes.
These instruments are based on the Sagnac effect, so called after the French
physicist Georges Sagnac. This effect states that two beams of light coming from
the same source and travelling along the same path but in opposite directions are
subject to a shift Du of phase, which depends on the angular velocity x of the
apparatus with respect to an immovable system of reference.
As shown in the following figure, a beam of light is split in two at a point P, and
the two resulting beams are made to follow the same path along the same circular
ring of radius r, but in opposite directions. When the circular ring, through which
light propagates, rotates clockwise at an angular velocity x 6¼ 0, then the beam of
light travelling clockwise (marked in blue) goes a bit further than the beam trav-
elling counterclockwise (marked in red).
Following Galloway [21], let ‘CW and ‘CCW be the lengths of the paths along,
respectively, the clockwise direction and the counterclockwise direction.
Let s = 2pr/c be the time taken by light to travel around the circular ring, where c
denotes the speed of light in vacuo.
We assume at first the two beams of light to propagate in vacuo. In this case, the
lengths ‘CW and ‘CCW can be expressed as follows
‘CW ¼ 2pr þ xrs

‘CCW ¼ 2pr xrs
This is because the light beam splitter moves through a length xrs = 2pxr2/c
during the propagation time s.
The difference in propagation times of the two beams of light is
‘CW ‘CCW 2xrs 4pr 2 x

Dt ¼ ¼ ¼
c c c2
In case of monochromatic laser light having frequency m (or wavelength k = c/m),

the difference of phase of the two beams of light is
8p2 r 2 mx 8p2 r 2 x
Du ¼ 2pmDt ¼ ¼
c2 ck
Now we assume the two beams of light to propagate through a homogeneous

dielectric medium, whose refractive index is n. In case of x = 0, the two beams
travel at a speed of c/n in opposite directions and take a time equal to ns to go
around the circular ring. In case of x 6¼ 0, the lengths ‘CW and ‘CCW of the optical
paths computed above for the two beams are to be modified as follows
8.11 Optical Gyroscopes 1101
2pnxr 2
‘ CW ¼ 2pr þ
c
2pnxr 2
‘ CCW ¼ 2pr
c
This is because the light beam splitter moves through a length xrns = 2pnxr2/c
during the propagation time ns. Armenise et al. [22] note that, in this case, the speed
of light is no longer the same for both of the signals. In particular, the speed vCW for
the clockwise-propagating signal is
c
vCW ¼ þ f xr
n
and the speed vCCW for the counterclockwise-propagating signal is

c
vCCW ¼ f xr
n
where f = 1 − 1/n2 is the Fresnel–Fizeau drag coefficient. The additive or sub-

tractive term fxr on the right-hand side of the two preceding equations is due to the
drag of light propagating in a medium which moves uniformly.
The two counter-propagating waves take different times to arrive at the point P.
The difference in time is
‘ CW ‘ CCW 2pr þ 2pnxr 2 =c 2pr 2pnxr 2 =c

Dt ¼ ¼
vCW vCCW c=n þ f xr c=n f xr
Since c2/n2 fx2r2, then the preceding equation can be written as follows
4pr 2 n2 xð1 f Þ 4pr 2 x

Dt ¼
c2 c2
This is because f = 1 − 1/n2. Therefore, Dt* = Dt. In other words, the shift of
phase Du due to the angular velocity x has the same value, independently of
whether the two beams of light propagate in vacuo or in a homogeneous medium of
refractive index n. The shift of phase is measured by means of a ring interferometry
setup, because the interference fringes are displaced from their position when the
apparatus is not rotating. The equation derived above
8p2 r 2 x
Du ¼
ck
holds in case of a circular ring of radius r. In case of a coil made of a fibre optic
cable having k turns, the same expression becomes
8p2 r 2 kx
Du ¼
ck
These concepts have been shown above on an intuitive level, apart from rela-
tivistic considerations. However, the Sagnac effect can also be explained within the
framework of general relativity. Details on the matter can be found, for example, in
[23]. The basic configuration of a fibre optic gyroscope is shown in the following
figure (redrawn from Ref. [24]).
A fibre optic gyroscope requires calibration, that is, determination of the indi-
cation corresponding to zero angular velocity (x = 0).
An optical gyroscope of another type is the ring laser gyroscope. It is a laser,
whose cavity forms a ring or a closed optical circuit, as shown in the following
figure, which is due to the courtesy of the United States Air Force [25].
8.11 Optical Gyroscopes 1103
Fibre optic gyroscopes and ring laser gyroscopes are both of them applications
of the Sagnac effect described above. In a ring laser gyroscope, two counter-rotating
laser beams travel along a closed circuit or ring, which is triangular or rectangular
and is usually filled with gas. A mirror is located at each corner of the ring for the
purpose of turning the laser beams through a precise angle, so as to create a path of
exact length.
The length of the laser cavity determines the wavelength k (or the frequency m)
of the light which is emitted, because the wavelength in each direction must result
in an integral number of waves. A detector or an output sensor is placed at one of
the corners. Instead of measuring the difference Dt in times of travel, the detector
measures the difference Dm in frequency, by using the Doppler effect. The laser
beam travelling in the same direction as the rotation of the platform travels through
a longer distance and has therefore a lower frequency m than that of the other laser
beam.
The difference Dm in frequency is proportional to the angular velocity x of the
vehicle on which the ring laser gyroscope is mounted, as follows
4Ax
Dm ¼
npk
([21]), where A is the area enclosed in the ring laser, p is the perimeter of the
optical ring, and n is the refractive index of the gas which fills the laser cavity.
Ring laser gyroscopes are more accurate than fibre optic gyroscopes, but are
subject to a phenomenon known as frequency lock-in at very low absolute values of
angular velocity x. In the following figure, the hatched line indicates the ideal
response of a ring laser gyroscope, whereas the red line indicates the difference Dm
in frequency effectively measured by the instrument as a function of x. This
phenomenon happens because a ring laser gyroscope measures rates of turn of a
vehicle by sensing differences in frequency.
When a rate of turn has very low values, then the corresponding difference in
frequency is also very small, and therefore the two frequencies tend to couple
together, because the two frequencies become almost identical, causing the
instrument to indicate a zero turning rate. This frequency coupling arises from back
scattering of the mirrors [26]. In order to overcome this problem, the ring laser
gyroscope is mechanically dithered, that is, rapidly moved through the lock-in zone
(also called dead band). By so doing, the time constant of the system keeps the
gyroscope from getting locked. This is obtained by means of a sinusoidal motion
about the input axis of the gyroscope [27], for the purpose of suppressing the dead
band.
8.12 Vibrating Structure Gyroscopes
The gyroscopes described here are also called Coriolis vibratory gyroscopes,
because they are based on the effect of the Coriolis force on a vibrating body.
As is well known, the Coriolis force is a fictitious force acting on a body of mass
m which moves in a rotating system of reference. This force has the following
expression
f C ¼ 2mx vr ¼ 2mx r
where x is the angular velocity vector, having constant magnitude and direction, of
the rotating system with respect to an inertial system, and vr
r` is the velocity
vector of the mass m with respect to the rotating system of reference.
The suspended body of mass m shown in the following figure has two degrees of
freedom, since it can only move along the axes x and y of a rotating system of
reference xyz, whose z-axis (not shown in the figure) is perpendicular to the plane of
the sheet, so that uz = ux uy.
The angular velocity vector x whose components along x, y, and z are
x ¼ 0ux þ 0uy þ xuz

8.12 Vibrating Structure Gyroscopes 1105
is directed along the z-axis. The oscillation along the x-axis (primary oscillating
mode) is driven by a force f, whose components along the axes x, y, and z are
f ¼ fx ux þ 0uy þ 0uz
The oscillation along the y-axis (secondary oscillating mode) is only due to the
rotation of the body about the z-axis. Two springs (of stiffness 1=2 kx each) and two
dashpots (of damping coefficient 1=2 cx each) are placed horizontally.
Likewise, two springs (of stiffness 1=2ky each) and two dashpots (of damping
coefficient 1=2 cy each) are placed vertically.
The equation of motion for the body described above, written in the rotating
system of reference and projected onto the axes x and y, can be written as follows
mxþ cx xþ kx x ¼ fx þ 2mxy

myþ cy yþ ky y ¼ 2mxx
This is because the Coriolis force fC has the following components along the
axes x, y, and z
f C ¼ 2mx vr ¼ 2mx r ¼ 2mxyux 2mxxuy þ 0uz
In most cases, the driving force f = fxux is a sinusoidal function of time, such that
fx ¼ Ax sinðxd tÞ
where Ax and xd are, respectively, the amplitude and the angular frequency of fx.
The angular frequency xd of the driving force is kept close to the natural angular
frequency xx = (kx/m)1/2 of the primary oscillator, so that x can be expressed as a
function of time as follows
x ¼ Ax sinðxd tÞ Ax sinðxx tÞ
The function x, differentiated with respect to time, yields
x ¼ Ax xx cosðxx tÞ
The preceding expression, substituted into the second of the two equations
written above (my`` + cyy` + kyy = −2mxx`), leads to

xy
y ¼ yþ x2y y ¼ 2Ax xxx cosðxx tÞ
Qy
where Qy = (mky)1/2/cy is called the quality factor of the secondary oscillating

mode, and xy = (ky/m)1/2 is the natural angular frequency of the secondary
oscillator. We search the solution of the preceding differential equation in the

following form
y ¼ Ay cosðxx t þ uy Þ
where Ay and uy are two constants of integration, whose values are to be deter-
mined, as will be shown below. By differentiating once and twice the preceding
function with respect to time, we obtain
y ¼ Ay xx sinðxx t þ uy Þ
y ¼ Ay x2x cosðxx t þ uy Þ
Substituting these expressions of y, y`, and y`` into the preceding differential
equation y`` + (xy/Qy)y` + x2y y = −2Axxxx cos(xxt) leads to

Ay xx 2 þ Ay xy 2 cos uy Ay xx xy =Qy sin uy cosðxx tÞ þ ðAy xx 2

Ay xy 2 Þ sin uy Ay xx xy =Qy cos uy sinðxx tÞ ¼ 2Ax xxx cosðxx tÞ
The preceding equation implies

Ay xx 2 þ Ay xy 2 cos uy Ay xx xy =Qy sin uy ¼ 2Ax xxx

ðAy xx 2 Ay xy 2 Þ sin uy Ay xx xy =Qy cos uy ¼ 0
This is a system of two algebraic equations for the two unknowns Ay and uy.
The second equation yields immediately
xx
tan uy ¼ x y
Qy x2x x2y
The first equation can be rewritten by using the preceding expression of tan uy
and the following trigonometric identities
tan uy
sin uy ¼ 12
1 þ tan2 uy
1
cos uy ¼ 12
1 þ tan2 uy
8.12 Vibrating Structure Gyroscopes 1107
This yields
2Ax xx
Ay ¼ " #1 x
2 x x 2 2
x y
x2x x2y þ
Qy
This equation shows that the amplitude Ay of the secondary oscillating mode is
proportional to the angular velocity x. Therefore, the preceding equation
y ¼ Ay cosðxx t þ uy Þ

8 2 39
2Ax xx x < x x =
xx t þ arctan4 5
x y
y¼" # 1 cos
2 x x 2 2 : Qy x2x x2y ;
x y
x2x x2y þ
Qy
In conclusion, the angular velocity x of a body can be measured, by means of a

Coriolis vibratory gyroscope, by:
• applying a periodic driving force, having an appropriate angular frequency, to
the body along the horizontal direction; and
• measuring the amplitude of the oscillation of that body along the vertical
direction.
Further information on gyroscopes of this type can be found in Refs. [28, 29].
8.13 Accelerometers
As has been shown in the preceding sections, gyroscopes provide information of

directional type concerning angles and their rates of change with time. The present
section deals with accelerometers, which provide information on forces acting on
objects. When the acceleration of an object is made available for each direction,
then the velocity of that object and the distance travelled by it can be obtained by
means of two successive integrations. Therefore, given the position of a space
vehicle at an initial time, its position at any other subsequent time can be deter-
mined by means of a process called dead reckoning.
Following the definition given in Ref. [30], an accelerometer is “a device which
uses a proof mass to provide an output that is a known function of acceleration.
Therefore, the accelerometer includes the electromechanical parts such as the proof
mass, type of restraint, pickoff, and other electronic parts required to provide the
output”. The output mentioned above is usually a voltage signal, which varies in
proportion to the acceleration applied to the proof mass, as will be shown below.
The following figure, due to the courtesy of NASA [30], shows the scheme of a
linear accelerometer based on a suspended mass.
Another type of linear accelerometer is the pendulous mass accelerometer, which

is shown in the following figure.
A pendulous mass accelerometer uses an internal proof mass on a cantilever to

sense the orientation of the case with respect to an inertial system of reference. Two
position sensors use electrostatic forces to centre the proof mass in a chamber. The
position of the proof mass is detected by either sensor, and a rebalance force is
applied by means of two coils to keep the mass centred. The acceleration signal is
derived from the rebalance force signal.
The acceleration applied to a space vehicle acts on the case of the accelerometer
carried on board that vehicle. When the case accelerates, the proof mass m con-
tained in the accelerometer tends to remain at rest (or move at a uniform velocity)
with respect to an inertial system of reference, in accordance with Newton’s first
law. Therefore, the proof mass moves with respect to the case of the accelerometer.
The restraint mechanism applies a force to the proof mass. This force accelerates the
proof mass and puts it in a condition of equilibrium with respect to the case. In this
condition, the force f exerted by the restraint mechanism on the proof mass is
8.13 Accelerometers 1109
proportional to the acceleration a of the proof mass (which acceleration is the same
as that of the case), in accordance with Newton’s second law, that is,
f ¼ ma
The accelerometer senses the mechanical acceleration applied to the proof mass
and generates, by means of a transducer, an electric signal which is proportional to
the force applied.
An angular accelerometer works in a manner similar to that described above. It
uses the moment of inertia I of a balanced proof mass to exert a torque M in
response to an angular acceleration a″ applied to it. This is also in accordance with
Newton’s second law, that is,
M ¼ Ia00
The accelerometers used for space applications are usually linear accelerometers.
However, angular accelerometers also exist and have been used in space vehicles.
The acceleration applied to a space vehicle is due to gravitational and
non-gravitational forces. An accelerometer responds to any force which generates a
motion between the proof mass and the case, independently of the nature of that
force. For example, for a space vehicle in free fall, the accelerometer mounted on it
has no input, and its output is only the accelerometer bias. For a space vehicle
subject only to thrust, the accelerometer output comprises the bias and the accel-
eration due to thrust. An accelerometer placed on the surface of the Earth is subject
to the gravitational force, and therefore its output comprises the bias and the
acceleration due to gravity. Finally, an accelerometer placed on the surface of the
Earth and moved from its rest position by a propulsive force is subject to a total
force which results from the vector sum of the gravitational force and the propulsive
force.
Linear accelerometers are meant to sense translational accelerations acting along
one axis. This axis is called the input axis of the accelerometer. A linear
accelerometer should be designed in such a way as to have the minimum possible
response to accelerations laying in a plane perpendicular to its input axis.
This instrument is basically a case containing a proof mass restrained between
two springs, a damper, and a pickoff. When the centre of the proof mass is displaced
from its rest position (corresponding to no external forces applied to the case), the
restraining springs generate a force which acts on the proof mass and accelerates it.
In case of a constant acceleration impressed, the centre of the proof mass reaches a
position of equilibrium which is placed at a distance from the rest position. In the
condition of equilibrium, the force exerted by the springs on the proof mass is
proportional to the acceleration acting on the case, in accordance with Newton’s
second law. As is well known, the restraint force produced by a spring attached to a
mass is proportional to the displacement of that mass. Therefore, the displacement
of the centre of the proof mass is proportional to the acceleration acting on the case.
An ideal accelerometer has only one degree of freedom along the input axis and is
insensitive to accelerations acting on a plane perpendicular to this axis. The

methods used for suspending the proof mass within the case reach this ideal goal
with various degrees of exactness.
The working principle of a linear accelerometer can be described mathematically
by considering the following figure (redrawn from Ref. [31]).
Let m, x, and d be, respectively, the proof mass, the displacement of this mass
from its rest position, and the displacement of the case from an inertially fixed
point. Let k and c be, respectively, the stiffness of the spring system and the
damping coefficient of the dashpot, which has the function of dissipating energy.
Following Widnall et al. [31], the total force fx applied to the proof mass m in the
horizontal direction is
fx ¼ mðd þ xÞ00
The force fx is counterbalanced by the reactions due to the spring system (−kx)
and the dashpot (−cx′), and therefore
fx ¼ kx cx0
In other words, the equation which governs the motion of the mass m is
mx00 þ cx0 þ kx ¼ md 00
where the term on the right-hand side is the driving force. The characteristics x and
f of a linear accelerometer result from the corresponding homogeneous equation
x00 þ 2fxx0 þ x2 x ¼ 0
where x = (k/m)1/2 is the natural angular frequency of the instrument in the absence
of damping, f = c/ccr is the damping ratio, and ccr = 2(km)1/2 is the critical
damping. According to [31], typical values of x and f are
8.13 Accelerometers 1111
x ¼ 103 rad=s
f ¼ 0:7
Assuming the driving acceleration −d″ of the case to be a sinusoidal function of

time, the distance d can be expressed as follows
d ¼ Ad sinðxd tÞ
where Ad and xd are, respectively, the amplitude and the angular frequency of the
driving acceleration. In this case, the preceding equation mx″ + cx′ + kx = −md″
mx00 þ cx0 þ kx ¼ mAd x2 d sinðxd tÞ
The steady-state solution of this non-homogeneous differential equation is a

sinusoidal function of time
x ¼ A sinðxd t uÞ
which has the same angular frequency xd as that of the driving acceleration, and an
amplitude A and a phase angle u, whose values are to be determined.
By using the method shown in Sect. 8.12, we find
" #
2fðxd =xÞ
u ¼ arctan
1 ðxd =xÞ2
Ad ðxd =xÞ2
A¼ 12
h i2
1 ðxd =xÞ2 þ ½2fðxd =xÞ 2
By multiplying both terms of the second of the two equations written above by
x2, there results
Ad x2d
Ax2 ¼ 12
h i2
1 ðxd =xÞ2 þ ½2fðxd =xÞ 2
which is the amplitude of the acceleration of the case, because
d 00 ¼ Ad x2 d sinðxd tÞ
where the minus sign means that the acceleration −d″ of the case is out of phase by
p radians with respect to the displacement d of the case itself.
When the angular frequency ratio xd/x is much smaller than unity, then
Ax2 Ad x2 d
An accelerometer having a small mass and a stiff spring has a high natural
angular frequency x = (k/m)1/2. Therefore, when the angular frequency xd of the
driving acceleration is much smaller than x, then the displacement of the mass m is
proportional to the acceleration of the case. In other words, by measuring the
displacement A of the mass m within the case and multiplying this displacement by
the square of the known value x of the natural angular frequency of the
accelerometer, we obtain the acceleration of the case.
Three accelerometers placed together in mutually perpendicular directions make
it possible to determine the acceleration in any direction of interest.
An electric signal is generated in an accelerometer by means of a displacement
transducer which may be of various types. Common types of transducers are
described in Refs. [32, 33]. They are:
• capacitive, in which a metal beam attached to the proof mass alters the distance
between two metal plates, and therefore a measurement of capacitance gives a
measurement of the acceleration impressed;
• piezoelectric, in which single-crystal quartz (SiO2) or other piezoelectric
material (such as barium titanate BaTiO3, lead titanate PbTiO3, …) attached to
the proof mass is compressed when the proof mass moves, and therefore gen-
erates a tiny voltage which is proportional to the acceleration impressed;
• piezoresistive, in which the proof mass is attached to a potentiometer, or vari-
able resistor, which causes an electric current to vary in proportion to the
acceleration impressed;
• Hall effect, in which the acceleration impressed to the proof mass is measured by
sensing tiny changes in the magnetic field around the accelerometer;
• magneto-resistive, in which the resistivity of a material changes in the presence
of a magnetic field; and
• heat transfer, in which changes in heat transfer due to the acceleration of the
proof mass are measured by the asymmetrical heat gradient.
8.14 The Stable Platform
A stable platform is a surface mounted on gimbals and containing gyroscopes and

accelerometers. It is an essential part of an inertial guidance system and is used to
maintain a desired orientation with respect to an inertial system of reference,
independently of the motion of the vehicle on which it is carried.
8.14 The Stable Platform 1113
A stable platform uses the property of gyroscopic precession (a torque about an

input axis, except the spin axis, produces an angular velocity about the orthogonal
output axis). For this purpose, three gyroscopes, each having one degree of free-
dom, are mounted on a stable platform along three mutually perpendicular direc-
tions. Three accelerometers are also mounted on the platform along three mutually
perpendicular directions, as shown in the following figure, due to the courtesy of
the United States Air Force [25], which illustrates an aircraft autopilot, where three
ring laser gyroscopes are used to sense changes in pitch, roll, and yaw. In addition,
there are three pendulous mass accelerometers to measure longitudinal, lateral, and
vertical motion.
A stable platform is mounted on gimbals, as shown in the following figure,

which is due to the courtesy of NASA-MIT [34]. This figure illustrates a scheme of
the three-degree-of-freedom servo-driven gimbals which isolate the gyroscopes and
accelerometers from the motion of the base.
If the platform were perfectly balanced and the bearings were frictionless, no
torque would act on the platform, and its orientation would be maintained inde-
pendently of the motion of the carrying vehicle. In practice, since unbalances and
friction cannot be eliminated completely, the platform is subject to perturbing
torques. The three gyroscopes sense such torques and counteract them by means of
a servo system, in order to obtain a torque-free platform.
In order to describe the working principle of a stable platform, let us consider the
simple scheme illustrated in following figure. This scheme shows a stable platform,
which can only rotate about a single axis. Let y be the rotation axis of the platform,
that is, the input axis. Let Uy be the angle through which the platform rotates about
the y-axis as a result of a perturbing torque My. The other two axes of the body-fixed
system xyz are such that the x-axis is the output axis, which is perpendicular to y,
and the z-axis, which is perpendicular to the xy-plane. The perturbing torque My,
which acts about the y-axis, causes the spin axis of the gyroscope to rotate about the
x-axis through an angle h. Consequently, the moment-of-momentum vector h also
rotates about the x-axis through the same angle h.
The balance of the moments acting about the y-axis requires the perturbing
torque My minus the inertia torque Jy U″y (also acting about the y-axis) to be equal
to the rate of change of the moment-of-momentum vector h.
The balance of the moments acting about the y-axis is expressed by the fol-
lowing equation
My Jy U00 y ¼ hh0
where Jy is the moment of inertia of the system comprising the platform and the
gyroscope with its gimbal suspension about the y-axis. The term on the right-hand
side of the preceding equation results from an approximation, which is justified,
because the value of the angle h is usually less than 1° [10].
This is because, as has been mentioned above, the perturbing torques acting on a
stable platform are counteracted by a servo system, which is not shown in the
simple scheme of preceding figure.
A better scheme of a stable platform having a single degree of freedom is shown
in the following figure, which is due to the courtesy of NASA [35].
As shown in this figure, information coming from the signal generator is used in
the servo motor to keep the platform oriented. The torque generator is used to
correct errors in the gyroscope and to calibrate it. Any rotation of the platform about
the input axis (IA in the figure) of the gyroscope is sensed by the signal generator,
amplified, and then used to drive the servo motor to return the platform to the initial
orientation [35].
The balance of the moments acting about the x-axis requires the total torque
acting about the x-axis to be equal to zero, according to the following equation
Ix h00 hU0 y ¼ 0
where Ix is the moment of inertia of the gyroscope and its gimbal suspension about
the x-axis.
We use again the notation f ðsÞ

Lf ðtÞ to indicate the Laplace transform of a
function f(t). By so doing, we have Lh0 ¼ shðsÞ, and then the two preceding
equations of motion can be written as follows
M y ðsÞ Jy s2 Uy ðsÞ ¼ hshðsÞ

Ix s2 hðsÞ ¼ hsUy ðsÞ
The second of these equations, solved for Uy ðsÞ, yields
Uy ðsÞ ¼ Ix shðsÞ=h
The preceding expression is substituted into the first of the two equations written
above, which is turn is solved for hðsÞ=M y ðsÞ. This yields
h
hðsÞ Jy I x
¼
M y ðsÞ h2
s s2 þ
Jy Ix
which is the transfer function between the output function hðsÞ and the input
perturbing torque M y ðsÞ.
The angular velocity h′ of the frame of the gyroscope with respect to the plat-
form is sensed by an electric pickoff, amplified, and then fed to a servo motor,
which generates a counter-torque Ms, that is, a torque acting in the direction
opposite to that of the perturbing torque My. Since the value of JyIx is much greater
than that of h2, then the preceding equation can be approximated as follows
hðsÞ h

M y ðsÞ J y xs
I 3
Therefore, the servo system of a stable platform with a single degree of freedom can
be represented by the block diagram illustrated in the following figure, where A(s) is
the transfer function of the electric pickoff, the amplifier, and the servo motor. When
the transfer function A(s) is known, then the dynamical behaviour of the stable platform
can be studied for stability.
Let us consider now a platform stabilised on three axes. This platform may be
considered as an assembly of three single-axis platforms like that described above,
but mounted on gimbals, so as to make a single stable unit. The analysis of a
three-axis stable platform is complex, because of the coupling of each rotation with
the other two, and also because of the necessity of resolving the pickoff signals,
since the axes of the platform are not aligned with the axes of the gimbals, as shown
in the following figure (redrawn from Ref. [10]), which illustrates the
moment-of-momentum vectors hx, hy, and hz of the x-, y-, and z-gyroscopes.
Let hx, hy, and hz be the outputs of the x-, y-, and z-gyroscopes in response to the
rotations Ux, Uy, and Uz of the platform about the respective input axes. By
inspection of the preceding figure, the pickoff signal of each gyroscope results
r x ¼ hx U y
ry ¼ hy þ Ux
rz ¼ hz Uy
When the gimbal axes of the gyroscopes are aligned with the platform axes, then
the counteracting torques generated by the pickoff signals are of the following form:
Msy ¼ Ay ðsÞ

ry ðsÞ ¼ Ay ðsÞ hy ðsÞ þ Ux ðsÞ
where Ay(s) is the transfer function of the y servo system. Therefore, the block
diagram shown above (relating to a stable platform with a single degree of freedom)
must be modified as shown in the following figure (redrawn from Ref. [10]), in
order to take account of the coupling term of the form Ay(s)Ux(s).
The block diagram of the platform stabilised on three axis comprises three
uncoupled circuits like that relating to a platform stabilised on a single axis (y), and
additional connexions corresponding to the coupling terms. These terms are
−Ax(s)Uy(s), Ay(s)Ux(s), and −Az(s)Uy(s), as shown in the preceding figure.
We assume that the outer gimbal axis, originally parallel to the x-axis of the
platform, is attached to the vehicle, as shown in the following figure.
In addition, we assume that the vehicle is stabilised about its roll axis, which
implies Ux = 0. In this case, the rotations of the vehicle in pitch (Uy) and yaw (Uz)
cause the gimbal axes to deviate from the platform axes. Therefore, the torques Mx,
My, and Mz acting on the platform must be resolved along the displaced gimbal
axes, where the servo motors act. Since the counter-torques are proportional to the
pickoff signals of the platform, then the proper torques acting about the new gimbal
axes result from resolution of the pickoff signals of the platform along the gimbal
axes. Remembering that the vehicle is stabilised about its roll axis, we first allow
the nose of the vehicle to go down through a pitch angle Uy. Let x*, y*, and z* be
the new gimbal axes.
The components of the pickoff signals of the platform along the gimbal axes are
rx ¼ rx cos Uy rz sin Uy
ry ¼ ry
rz ¼ rx sin Uy rz cos Uy
Then, we allow a rotation Uz in yaw about the z*-axis and resolve rx* along the
pitch gimbal axis and the new roll axis, as follows
rroll ¼ rx sec Uz
rpitch ¼ rx tan Uz
Therefore, the resulting pickoff signals about the new gimbal axes of roll, pitch,
and yaw due to the two rotations Uy and Uz are

rroll ¼ rx cos Uy rz sin Uy sec Uz

rpitch ¼ rx cos Uy rz sin Uy tan Uz þ ry
ryaw ¼ rz ¼ rx sin Uy rz cos Uy
The equations written above can also be expressed in matrix notation as follows
2 3 2 32 3
rroll cos Uy sec Uz 0 sin Uy sec Uz rx
4 rpitch 5 ¼ 4 cos Uy tan Uz 1 sin Uy tan Uz 54 ry 5
ryaw sin Uy 0 cos Uy rz
The resolver performs the function of resolving the platform pickoff signals rx,
ry, and rz to the components rroll, rpitch, and ryaw along the displaced gimbal axes
where the roll, pitch, and yaw servo motors are placed.
8.15 Inertial Navigation
By navigation we mean “the process of planning, recording, and controlling the

movement of a craft or vehicle from one place to another” (Ref. 12, p. 799). There
are several manners in which this object can be attained. Reference [36] mentions in
particular those indicated below:
• piloting, which consists in using landmarks to determine the current position of
a vehicle;
• dead reckoning, which consists in determining the current position of a vehicle
by using another position of the same vehicle determined previously and
information on its heading and speed over an elapsed time;
• celestial navigation, which consists in determining the current position of a
vehicle by using time and angular measurements taken between the local ver-
tical (or the visible horizon) and a celestial body of reference (the Sun, the
Moon, a planet or a star);
• radio navigation, which consists in determining the current position of a vehicle
by using either radio beacons placed in known positions on the surface of the
Earth or a Global Navigation Satellite System (GNSS); and
• inertial navigation, which consists in determining the current position of a
vehicle by using another position of the same vehicle determined previously and
then an on-board computer which integrates measurements (relating to an
inertial system) taken by accelerometers and gyroscopes, without the need to
rely on external references.
These forms of navigation may used either individually or in combination with
one another. The subject of the present section is inertial navigation, which is a
technique using the output signals from accelerometers and gyroscopes to track the
position and orientation of a vehicle with respect to a known starting point, ori-
entation, and velocity [37]. With reference to the following figure, which is due to
the courtesy of NASA [38], an inertial navigation system (INS), sometimes called
an inertial measurement unit (IMU), is self-contained within the carrying vehicle
and can also be used in the absence of radio navigation installations or a global
positioning system (GPS).
This system (or unit) contains three orthogonal rate gyroscopes and three
orthogonal accelerometers, which measure, respectively, angular velocities and
linear accelerations. The signals coming from these devices are processed by a
computer in order to determine the position and orientation of the carrying vehicle.
Such systems are used not only in spacecraft, but also in aircraft, missiles, sub-
marines, and ships. The present discussion is limited to those used in spacecraft.
Inertial navigation systems may be classified into two categories, depending on
the system of reference in which they operate. They are:
• systems mounted on stable platforms; and
• systems mounted directly on vehicles (strap-down).
Stable-platform systems are based on the stable member on which the inertial
sensors (accelerometers and gyroscopes) are mounted. This stable member is iso-
lated by means of gimbals from any external rotational motion, so as to be con-
stantly aligned with an inertial system of reference.
Three gimbals would technically suffice for this purpose. However, most
stable-platform systems use four gimbals in order to avoid a phenomenon called
gimbal lock, which occurs when two of the three gimbal axes become aligned.
Since each gyroscope operates about a single axis, then in case of a gimbal lock
there are only two independent axes instead of three, and therefore the capability of
isolation on three axes is inhibited [39]. The fourth (redundant) gimbal has a motor
which keeps it always oriented away from the other gimbals. By so doing, there are
always three independent axes, as shown in the following figure, which is due to the
8.15 Inertial Navigation 1123
The three gyroscopes mounted on the stable member measure angular rates, and
the gimbal driving system (consisting in torque motors) uses the angular rate
information to rotate the gimbals, in order to cancel out the rotations sensed by the
gyroscopes. In this manner, the gyroscopes and accelerometers mounted on the
stable member are inertially stabilised from the motion of the vehicle, and the stable
member represents physically an inertial system of reference (Ref. [41]).
The information on the attitude of the vehicle is obtained from the gimbal angles
of the stable member. The angles between adjacent gimbals can be read by means
of pickoffs. The signals coming from the three accelerometers are integrated two
times in order to calculate the position of the stable member. For this purpose, it is
necessary to subtract previously the term due to the gravitational acceleration from
the signal relating to the acceleration along the vertical axis, as shown in the
following scheme, which is redrawn from Ref. [41].
Systems mounted on stable platform provide a good dynamic environment for

inertial sensors. This holds in particular in conditions of high angular oscillatory
rates. The experience gained so far with such systems yields operation conditions
characterised by very low navigational errors, which are mainly due to uncertainties
in the determination of the gravity field to which the vehicle is subject [41].
In strap-down systems (illustrated in the following scheme, which is redrawn
from Ref. [41]), all the inertial sensors are mounted, either directly or through
vibration isolators, on a platform which changes orientation together with the
vehicle. Consequently, the quantities (accelerations and angular rates) measured by
such sensors relate to the body-fixed system of reference.
In order to use an inertial reference system, a computer executes a procedure
based on a transformation matrix (denoted with RT in the following scheme)
between the body-fixed system and the inertial system. This procedure processes
the signals coming from the gyroscopes, as the vehicle moves and changes its
orientation. Successively, the signals coming from the accelerometers are processed
in order to be transformed from the body-fixed system to the inertial system. The
inertial sensors which are mounted on strap-down systems must have a larger
dynamic range than that of the sensors mounted on inertial-platform systems,
because the former are subject to the forces which act on the vehicle.
On the other hand, the absence of gimbals in the latter reduce their size and cost.
The increased computational requirements typical of strap-down systems tend to
become of minor importance due the development of computer technology.
Let us consider the transformation of the signals coming from the accelerometers
from the body-fixed system to the inertial system.
Following Titterton and Weston [42], let f b and f i be the specific force vector,
that is, the acceleration vector, expressed, respectively, in the body-fixed system
and in the inertial system. Here, the superscripts b and i indicate, respectively, the
body-fixed system and the inertial system. With reference to the preceding figure,
the components of fb are the signals coming from the three accelerometers.
The specific force vector fi in the inertial system results from pre-multiplying fb
by the rotation matrix RT, as follows
f i ¼ RT f b
where the rotation matrix RT defines the orientation of the body-fixed system with
respect to the inertial system. This matrix results from the angular rates measured by
the three gyroscopes, as follows

RT ¼ RT Xb ib
where Xb ib is the following skew-symmetric matrix

2 3
0 xz xy
Xbib ¼ 4 xz 0 xx 5
xy xx 0
The matrix Xb ib is formed by using the three components of the vector

2 3
xx
xbib ¼ 4 xy 5
xz
These three components represent the turn rates of the body with respect to the
inertial system as measured by the three gyroscopes. The rotation matrix RT,
needed for a transformation of the acceleration vector fb from the body-fixed system
to the inertial system, is the transpose of the rotation matrix R, which will be shown
below.
Let us consider a spacecraft rotating about its centre of mass O, and an inertial
system of reference XYZ. Let x, y, and z be the principal axes of inertia having their
origin in the centre of mass of the spacecraft, as shown for the Space Shuttle in the
following figure, due to the courtesy of NASA [43].
The translational motion of O with respect to the inertial system XYZ results
from a double integration of the signals coming from three orthogonal
accelerometers.
As to the rotational motion of the spacecraft, the angles Ux (roll), Uy (pitch), and
Uz (yaw) indicate the orientation of the body-fixed system xyz with respect to the
It is to be noted that the yaw, pitch, and roll angles are not the classical Eulerian
angles w, h, and / defined in Sect. 1.8. They differ from the classical Eulerian
angles because, for the yaw, pitch, and roll angles, only one rotation takes place
about each axis (3-2-1 sequence); whereas, for the classical Eulerian angles, two
rotations take place about one axis (3-1-3 sequence).
The rotations from the inertial axes XYZ to the body-fixed axes xyz are as
follows:
• first, a rotation Uz about the Z axis;

• then, a rotation Uy about the new Y axis; and
• finally, a rotation Ux about the final position of the X axis.
Therefore, the sequence of rotations indicated above is called the 3-2-1
sequence.
Let R1(Ux), R2(Uy), and R3(Uz) be the elementary matrices which define the
three rotations Ux, Uy, and Uz about to the axes, respectively, X, Y, and Z. These
matrices have the following components
2 3 2 3
1 0 0 cos Uy 0 sin Uy
6 7 6 7
R1 ¼ 4 0 cos Ux sin Ux 5 R2 ¼ 4 0 1 0 5
0 sin Ux cos Ux sin Uy 0 cos Uy
2 3
cos Uz sin Uz 0
6 7
R3 ¼ 4 sin Uz cos Uz 0 5
0 0 1
The rotation matrix R(Ux, Uy, Uz) = R1 R2 R3, which defines the transformation
2 3 2 3
x X
4 y 5 ¼ R Ux ; Uy ; Uz 4 Y 5
z Z
from the inertial co-ordinates XYZ to the body-fixed co-ordinates xyz, is

2 3
cUy cUz cUy sUz sUy
R ¼ 4 sUx sUy cUz cUx sUz sUx sUy sUz þ cUx cUz cUy sUx 5
cUx sUy cUz þ sUx sUz cUx sUy sUz sUx cUz cUy cUx
where ca and sa stand for, respectively, cos a and sin a (a being any angle).
A complete analysis of the matter has been made by Diebel [44].
Consequently, the rotation matrix RT, which defines the transformation from the
body-fixed co-ordinates xyz to the inertial co-ordinates XYZ, is
2 3
cUy cUz sUx sUy cUz cUx sUz cUx sUy cUz þ sUx sUz
RT ¼ 4 cUy sUz sUx sUy sUz þ cUx cUz cUx sUy sUz sUx cUz 5
sUy cUy sUx cUy cUx
When the roll, pitch, and yaw angles are about zero, then their sines can be
approximated by the respective angles expressed in radians (sin Ux Ux, sin Uy Uy,
and sin Uz Uz), and their cosines can also be approximated by unity (cos Ux 1,
cos Uy 1, and cos Ux 1). In this case, the products of small angles are also
negligible, and the rotation matrix RT written above is approximately
2 3
1 Uz Uy
RT 4 Uz 1 Ux 5
Uy Ux 1
This approximation is used to represent the small change in attitude which

occurs between a measurement and the successive update of this measurement in a
real-time computation of the body attitude. As has been shown above, the angles
Ux, Uy, and Uz are the gimbal angles, and their rates of change with time U`x, U`y,
and U`z are the gimbal rates. As will be shown in Sect. 9.1, the gimbal rates are
related to the body rates xx, xy, and xz as follows
2 3 2 3 2 3 2 3
xx 0 0 Ux
4 xy 5 ¼ R1 ðUx ÞR2 Uy R3 ðUz Þ4 0 5 þ R1 ðUx ÞR2 Uy 4 Uy 5 þ R1 ðUx Þ4 0 5
xz Uz 0 0
This yields
xx ¼ Uz sin Uy þ Ux

xy ¼ Uz cos Uy sin Ux þ Uy cos Ux
xz ¼ Uz cos Uy cos Ux Uy sin Ux
The equations written above can also be solved for U`x, U`y, and U`z. This yields
Ux ¼ xx þ ðxy sin Ux þ xz cos Ux Þ tan Uy

Uy ¼ xy cos Ux xz sin Ux
Uz ¼ ðxy sin Ux þ xz cos Ux Þ= cos Uy
These equations may be solved in a strap-down system to update the Euler

equations with respect to the chosen reference system. However, the solutions of
the first and the third of these equations become indeterminate when Uy = ±p/2.
The rotation matrix RT, which is used for the transformation fi = RT fb, changes
with time, and its rate of change (RT)` is by definition
T DRT RT ðt þ DtÞ RT ðtÞ

R ¼ lim ¼ lim
Dt!0 Dt Dt!0 Dt
where RT(t) and RT(t + Dt) are the rotation matrices at the times, respectively, t and
t + Dt. The matrix RT(t + Dt) can be expressed as the product of two matrices, as
follows
RT ðt þ DtÞ ¼ RT ðtÞAðtÞ
where A(t) is the rotation matrix which relates the body system at the time t to the
body system at the time t + Dt. In case of small angles of rotation DUx, DUy, and
DUz, the rotation matrix A(t) can be expressed as follows
AðtÞ ¼ I þ DW
where I is the 3 3 identity matrix, and DW is the following 3 3 matrix

2 3
0 DUz DUy
DW ¼ 4 DUz 0 DUx 5
DUy DUx 0
where, again, DUx, DUy, and DUz are the small angles of rotation through which the
body-fixed system has rotated during the interval Dt about its axes of, respectively,
roll, pitch, and yaw. When Dt approaches zero, then the small-angle approximation
is valid.
Substituting RT(t + Dt) = RT(t) A(t) in the equation
T DRT RT ðt þ DtÞ RT ðtÞ

R ¼ lim ¼ lim
Dt!0 Dt Dt!0 Dt
yields
DW
RT ¼ RT lim
Dt!0 Dt
As the time interval Dt tends to zero, the difference quotient DW/Dt tends to the
skew-symmetric matrix Xb ib formed by using the three components of the vector
2 3
xx
xbib ¼ 4 xy 5
xz
which vector in turn represents the rotation rates of the body-fixed system with
respect to the inertial system expressed in the body-fixed axes x, y, and z, that is,
DW
Xiib ¼ lim
Dt!0 Dt
Therefore

RT ¼ RT Xiib
The preceding differential equation in matrix form may be solved by means of a

computer in a strap-down inertial navigation system to keep track of the attitude of
the vehicle with respect to the inertial system.
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Chapter 9
Attitude Stabilisation and Control of Earth
Satellites
9.1 Attitude of Earth Satellites
The translational motion of the centre of mass O of a rigid body of constant mass
m with respect to an inertial system XYZ is governed by the equation
f ¼ mvO þ x mvO
The rotational motion of the same body about its centre of mass O or about a
fixed point with respect to XYZ is governed by the equation
M O ¼ hO þ x hO
In many practical cases, these equations are not coupled and can be solved
independently of each other. In such cases, the attitude of a rigid body with respect
to XYZ depends only on the rotations of this body about its centre of mass. In other
words, these rotations are independent of the translation of the centre of mass of a
rigid body with respect to XYZ, and therefore the two motions, rotational and
translational, of this body can be considered separately.
An artificial satellite revolving about the Earth and having no moving parts can
be considered in many cases as a rigid body, whose rotational motion is governed
by the Euler equations, which follow from the second of the two equations written
above, that is, from
M O ¼ hO þ x hO
as has been shown in Sect. 7.6. The Euler equations are rewritten below for
convenience:

https://doi.org/10.1007/978-3-319-62220-0_9
1134 9 Attitude Stabilisation and Control of Earth Satellites

MOx ¼ Ix xx Iy Iz xy xz
MOy ¼ Iy xy ðIz Ix Þxz xx

MOz ¼ Iz xz Ix Iy xx xy
where MOx, MOy, and MOz are the components along the axes x, y, and z of the
external moment MO acting on the satellite, Ix, Iy, and Iz are the moments of inertia
of the satellite about its principal axes, O is the centre of mass of the satellite, xx,
xy, and xz are the components along x, y, and z of the angular velocity vector x
relative to the inertial space, and x, y, and z are the principal axes of inertia, whose
origin is the centre of mass O of the satellite, as shown in the following figure. The
attitude of the satellite, that is, the orientation of its principal axes of inertia x, y, and
z with respect to a chosen system of orthogonal axes xref, yref, and zref, is indicated
here by means of the angle of yaw Uz, the angle of pitch Uy, and the angle of roll Ux
(see Sect. 8.15).
The following figure shows these two sets of orthogonal axes and the angles of
roll (Ux), pitch (Uy), and yaw (Uz).
9.1 Attitude of Earth Satellites 1135
The sequence of rotations from the reference axes xref, yref, and zref to the
principal axes of inertia x, y, and z is the 3-2-1 sequence, which has been shown in
Sect. 8.15, and is given again below for convenience:
• first, a rotation Uz about the zref axis;
• then, a rotation Uy about the intermediate position of the yref axis; and
• finally, a rotation Ux about the final position of the xref axis.
We consider here not only the inertial system XYZ, the principal system of inertia
xyz, and the desired attitude system xrefyrefzref, but also the orbital-plane system
xoyozo, (subscript o) which is related to the revolution of the satellite in a circular
orbit around the Earth, and is such that its origin is the centre of mass O of the
satellite, its xo-axis coincides with the velocity vector of the centre of mass of the
satellite, its zo-axis is directed from the centre of mass of the satellite to the centre of
mass G of the Earth, and its yo-axis is perpendicular to the orbital plane of the
satellite, so that uyo = uzo uxo, where uxo, uyo, and uzo are the unit vectors of the
axes, respectively, xo, yo, and zo, as shown in the following figure.
Let R1(Ux), R2(Uy), and R3(Uz) be the elementary matrices which define the
three rotations Ux, Uy, and Uz about to the axes of, respectively, roll, pitch, and yaw.
As shown in Sect. 8.15, these matrices have the following components
2 3 2 3
1 0 0 cos Uy 0 sin Uy
6 7 6 7
R1 4 0 cos Uxsin Ux 5 R2 4 0 1 0 5
0 sin Ux cos Ux sin Uy 0 cos Uy
2 3
cos Uz sin Uz 0
6 7
R3 4 sin Uz cos Uz 0 5
0 0 1
The rotation matrix R(Ux, Uy, Uz) = R1R2R3, which defines the transformation
2 3 2 3
x X
4 y 5 ¼ R Ux ; Uy ; Uz 4 Y 5
z Z
from the inertial co-ordinates XYZ to the body-fixed co-ordinates xyz, is

2 3
cUy cUz cUy sUz sUy
R 4 sUx sUy cUz cUx sUz sUx sUy sUz þ cUx cUz cUy sUx 5
cUx sUy cUz þ sUx sUz cUx sUy sUz sUx cUz cUy cUx
where ca and sa stand for, respectively, cos a and sin a (a being any angle).
By using the transformation indicated above, the relations between the angular
velocities xx, xy, and xz of the satellite about its own principal axes of inertia x, y,
and z and the rates of change U`x, U`y, and U`z of the roll, pitch, and yaw angles
with respect to the inertial axes X, Y, and Z are given by
2 3 2 3 2 3 2 3
xx 0 0 Ux
4 xy 5 ¼ R1 ðUx ÞR2 Uy R3 ðUz Þ4 0 5 þ R1 ðUx ÞR2 Uy 4 Uy 5 þ R1 ðUx Þ4 0 5
xz Uz 0 0
in accordance with Ref. [27]. By executing these operations, there results
xx ¼ Uz sin Uy þ Ux

xy ¼ Uz cos Uy sin Ux þ Uy cos Ux
xz ¼ Uz cos Uy cos Ux Uy sin Ux
The equations written above can also be solved for U`x, U`y, and U`z. This yields
Ux ¼ xx þ ðxy sin Ux þ xz cos Ux Þ tan Uy

Uy ¼ xy cos Ux xz sin Ux
Uz ¼ ðxy sin Ux þ xz cos Ux Þ sec Uy
When the angles Ux, Uy, and Uz are small, sin a tan a a, and cos a 1,
where a is any of such angles. In this case, the preceding equations become
xx Uz Uy þ Ux Ux xx þ xz Uy

xy Uz Ux þ Uy Uy xy xz Ux
xz Uz Uy Ux Uz xy Ux þ xz
Let us consider now the orbital axes xo, yo, and zo (subscript o) illustrated in the
preceding figure. These axes rotate in space at a constant angular velocity
xo ¼ 0uxo xo uyo þ 0uzo
with respect to the inertial system XYZ, that is, the angular velocity vector xo is
aligned with the yo-axis but points in the opposite direction, as shown in the
preceding figure. The relations between the angular velocities xx, xy, and xz of the
satellite about its own principal axes of inertia x, y, and z and the rates of change
U`x, U`y, and U`z of the roll, pitch, and yaw angles with respect to the orbital axes
xo, yo, and zo are given by
2 3 2 3 2 3
xx Uz sin Uy þ Ux 0
4 xy 5 ¼ 4 Uz cos Uy sin Ux þ Uy cos Ux 5 þ R1 ðUx ÞR2 Uy R3 ðUz Þ4 xo 5
xz Uz cos Uy cos Ux Uy sin Ux 0
This yields
xx ¼ Uz sin Uy þ Ux xo cos Uy sin Uz

xy ¼ Uz cos Uy sin Ux þ Uy cos Ux xo ðsin Ux sin Uy sin Uz þ cos Ux cos Uz Þ
xz ¼ Uz cos Uy cos Ux Uy sin Ux xo ðcos Ux sin Uy sin Uz sin Ux cos Uz Þ
These equations can also be solved for U`x, U`y, and U`z. This yields
Ux ¼ xx þ ðxy sin Ux þ xz cos Ux Þ tan Uy þ xo sin Uz sec Uy

Uy ¼ xy cos Ux xz sin Ux þ xo cos Uz
Uz ¼ ðxy sin Ux þ xz cos Ux Þ sec Uy þ xo sin Uz tan Uy
in accordance with Ref. [11]. Again, when the angles Ux, Uy, and Uz are small, then
the preceding equations become
xx Uz Uy þ Ux xo Uz Ux xx þ xz Uy þ xo Uz

xy Uz Ux þ Uy xo Uy xy xz Ux þ xo
xz Uz Uy Ux þ xo Ux Uz xy Ux þ xz
The rotational motion of the satellite is governed by the Euler equations


Therefore, the attitude of the satellite at a given time t depends on:

• the angles of roll (Ux), pitch (Uy), and yaw (Uz) and the components xx, xy, and
xz of the angular velocity vector x of the satellite at a time t0 < t, for example,
at the time of the injection of the satellite into orbit; and
• the components MOx, MOy, and MOz of the total external moment vector MO
acting on the satellite during the time interval going from t0 to t.
According to the fundamental equation
M O ¼ hO þ x hO
the vector MO is equal to the change of the moment of momentum with time.
If no external moments were acting on a spacecraft, its initial rotational motion
would be maintained indefinitely. In practice, several types of external moments
may act on a spacecraft, which alter considerably this motion. According to Gerlach
[11], the principal causes of the perturbing moments which act on a spacecraft are
the following:
• aerodynamic forces;
• electromagnetic induction;
• solar radiation pressure;
• gravity gradient; and
• micrometeorites.
Sidi [27] also considers reaction moments produced by particles expelled from
the spacecraft, such as ions or molecules of exhaust gases. A brief description of
these causes will be given in the following sections.
9.2 Moments Due to Aerodynamic Forces
These moments result from the impingement of atmospheric gas particles on the
surfaces of a spacecraft. The aerodynamic forces of drag and lift acting on a satellite
revolving around the Earth have been considered in detail in Sect. 3.19. For the
purposes of the present section, we rewrite below the equation which expresses the
magnitude of the aerodynamic drag force

1 2
jf D j ¼ CD pA ¼ CD qjvrel j A
2
where CD is a dimensionless quantity called drag coefficient, p = ½q|vrel|2 is the

dynamic pressure on the moving vehicle, q is the atmospheric density at the height
of interest, |vrel| is the magnitude of the velocity of the vehicle with respect to the
atmosphere, and A is the projected area of the vehicle normal to the relative velocity
vector vrel. This vector, in turn, is the velocity of the incident air stream relative to
the vehicle. Assuming the atmosphere as rotating with the Earth as a rigid body, the
relative velocity vector can be expressed as follows
vrel ¼ v vatm
where v is the velocity of the satellite with respect to an inertial reference system,
and vatm is the velocity vector of the Earth atmosphere with respect to the same
reference system. The drag force vector, fD, is directed oppositely to the relative
velocity vector, and therefore the drag force is fD = −½CDqA|vrel|vrel.
The magnitude of the total aerodynamic moment acting on a satellite results
from the product
jf D j‘
where |fD| is the magnitude of the aerodynamic drag force fD expressed by the
equation written above, and ‘ is the moment arm. When the magnitude of the drag
force acting on a satellite and the position of the centre of pressure are known, then
the moment of this force with respect to the centre of mass can be computed, as has
been shown above.
The dynamic pressure for a satellite revolving around the Earth depends on its
altitude over the Earth surface and its relative velocity. The drag coefficient, CD, can
be evaluated with sufficient approximation in many practical cases under the fol-
lowing conditions: (a) the mean free path of the atmospheric particles is much
greater than the size of the satellite; and (b) the mean speed of the atmospheric
particles in the free stream is much smaller than the satellite speed. Thomson [50]
gives data on the dynamic pressure acting on a satellite in a circular orbit around the
Earth as a function of altitude. These data show that aerodynamic torques decrease
rapidly with altitude. In practical cases, according to Gerlach [11], aerodynamic
torques are of importance at altitudes up to about 800 km and may be used for
stabilisation purposes at altitudes up to about 500 km.
At altitudes usually considered suitable for aerodynamic stabilisation, between
120 and 600 km, the aerodynamic forces provide negligible damping, and therefore
some means of damping must be provided, such as eddy currents, viscous fluid
damping, or gyroscopes [9]. Wiggins [55] and successively Gerlach [11] have
calculated several aerodynamic moments acting on a hypothetical satellite of
cylindrical shape measuring 9.15 m in length and 1.51 m in diameter, subject to an
incident stream perpendicular to its longitudinal axis (z), as shown in the following
figure.
These calculations are based on the theory of the free flow of the atmospheric
particles, as indicated above. The centre of mass (O) of the satellite has been
assumed to be placed at a distance of 0.305 m from the geometric centre, which is
also the centre of pressure (CP). The values of the aerodynamic moments (roughly
0.1 N m at an altitude of 200 km, and 0.0001 N m at an altitude of 600 km) have
been compared with those relating to perturbing moments due to other causes, and
the complete results are given in Refs. [11, 55, 41]. The values indicated above vary
as a function of the solar radiation level (day or night, and solar cycle) in an
increasing measure as the height increases.
9.3 Moments Due to Electromagnetic Induction 1141
9.3 Moments Due to Electromagnetic Induction
These moments are induced by the magnetic field of the Earth, when a magnetic
field is present in a satellite. A magnetic field may be present in a satellite for
several reasons:
• permanent magnetism in the satellite;
• eddy currents induced in a rotating satellite or in its rotating parts; and
• currents flowing through a coil in the satellite, as is the case even when this
current has the sole purpose of producing a controllable torque on the satellite.
Gerlach [11] has found that, at altitudes greater than 150 km, the magnetic field of
the Earth can be approximated with sufficient accuracy by the field of a magnetic dipole
placed at the centre of the Earth. The axis of this dipole forms an angle of 18° with the
spin axis of the Earth. The components Ha (directed along the axis of the magnetic
dipole) and Hn (directed normally to the axis of the magnetic dipole) of the magnetic
field vector H of the Earth have been expressed in oersted by Gerlach as follows
1 3cos2 d sinð2dÞ
Ha ¼ 0:308 3 Hn ¼ 0:461 3
r r
rE rE
where rE is the radius of the Earth, and r and d are the polar co-ordinates, whose
origins are, respectively, the centre of the Earth and the axis of the magnetic dipole,
of the point P in which the two components Ha and Hn of the vector H are com-
puted, as shown in the following figure.
Let us suppose that a coil, having n windings, carries a current whose intensity is
i amperes. Let S (in cm2) and B (in gauss) be, respectively, the area enclosed by the
windings and the magnitude of the magnetic flux density vector B. The magnetic flux
density vector B is related to the magnetic field vector H by the following equation
B ¼ lH
where the value of the permeability l, in vacuo, is l0 = 1 gauss/oersted.

In case of the plane of the windings being parallel to the direction of the vector
B, the magnitude of the torque, in dyne centimetre, acting on the coil is
S
M¼ inB
10
As an example, for the cylindrical satellite described above and for a current of
1 A flowing in a coil consisting of only one winding around the largest cross
section of the satellite, Wiggins [55] has computed the resulting magnetic torque,
whose approximate value ranges from 0.0002 to 0.001 N m in the interval 200–
600 km of altitude. The complete results are given in Refs. [11, 55, 41]. Torques
having these values can be very useful for the purpose of controlling the attitude of
a satellite, and can easily be increased by using coils carrying a current of greater
intensity or having a larger number of windings.
However, no torques due to electromagnetic induction can be generated about an
axis parallel to the local magnetic field. In addition, since the two components Ha
and Hn of the magnetic field vector H are proportional to 1/r3, then the magnitude
of these torques decreases rapidly as the altitude increases. Gerlach [11] has
computed that, in case of satellites orbiting at geosynchronous altitude, the mag-
nitude of a torque due to magnetic induction is 1/300 of its value near the surface of
the Earth. Angular motions of satellites can also be damped by means of eddy
currents which induce magnetic moments. Conversely, in case of spin-stabilised
satellites, it may be desirable to counterbalance magnetic moments induced by eddy
currents, for the purpose of eliminating the decay of the spin rate of such satellites.
9.4 Moments Due to Solar Radiation Pressure
As has been shown in Sect. 3.16, the amount S0 of radiant energy emitted by the
Sun over all the frequencies, which falls each second on the unit area exposed
normally to the Sun rays at the mean Sun-Earth distance in the absence of the Earth
atmosphere, is called total solar irradiance. Kopp and Lean [10] have found recently
the value of 1360.8 ± 0.5 W/m2 for the total solar irradiance during the 2008 solar
minimum period. The magnitude, fS, of the force exerted by the solar radiation on a
surface A perpendicular to the direction of radiation is expressed in dyne by Gerlach
[11] as follows
9.4 Moments Due to Solar Radiation Pressure 1143
fS ¼ p0 ð1 þ RÞA
where A is expressed in cm2, p0 = S0/c is the ratio of the solar irradiance to the
speed of light in vacuo (c = 299792458 m/s), and R is a coefficient of reflection of
the surface. The value of p0 results from
p0 ¼ 1360:8=299792458 ¼ 4:539 106 N=m2 ¼ 4:539 105 dyne=cm2
which is not very far from the value (4.3 10−5 dyne/cm2) indicated by Gerlach.
A surface absorbing completely the solar radiation has a value of R equal to zero,
whereas this value is equal to unity for a perfect reflector. A transparent surface has
a value of R between 0 and −1. When the direction of the incoming radiation forms
an angle a with the normal to the surface stricken, then the magnitude of the force
exerted on this surface is, in case of perfect reflection, proportional to cos2a. The
moment of the force due to solar radiation about the centre of mass of the satellite
can be computed when the distance between the centre of pressure and the centre of
mass is known.
The magnitude of the torque due to solar radiation is nearly independent of the
altitude of the satellite over the surface of the Earth. At geosynchronous or higher
altitudes, this torque may be the principal term of the total torque acting on a
satellite. When a satellite is properly designed, the torque due to solar radiation can
be used to stabilise the satellite in attitude, provided that its altitude is not lower
than 500 km. In particular, Sohn [48] has shown that a space vehicle can be
stabilised against perturbing torques by attaching a suitable weathervane-type tail
surface to the vehicle. Acord and Nicklas [16] have also shown that solar radiation
pressure can be used for torque control in interplanetary spacecraft. Conversely, in
order to keep this type of perturbing torque to a low value, a satellite can be
designed in such a way as to have a symmetrical shape with respect to its centre of
mass [11].
For the purpose of determining the magnitude of this torque, Gerlach [11] and
Wiggins [55] have used the cylindrical model of satellite described above. The
results found by them indicate that, at altitudes less than 2000 km, the torque due to
solar radiation pressure is on the average twenty times as small as the torque due to
electromagnetic induction considered above. By contrast, at geosynchronous alti-
tude, the torque due to solar radiation pressure has been found to be fifteen times as
great as the torque due to electromagnetic induction.
9.5 Moments Due to Gravity Gradient
These moments arise because the gravitational force of the Earth is not exactly
equal for all the masses which make up the total mass of a satellite.
In particular, in case of a prolate cylindrical satellite like that described in

Sect. 9.2, the different values of the gravitational force acting on the various masses
generate a torque about the axis of symmetry of the satellite. This is because such
masses have different distances from the centre of mass G of the Earth. The fol-
lowing figure shows a satellite revolving around the Earth and having its centre of
mass in the point O.
Let rO GO be the position vector of the centre of mass of the satellite in the
geocentric equatorial system XYZ. Let dm and r GP be, respectively, an
infinitesimal mass placed in a point P other than O and its position vector in the
geocentric equatorial system. Let q OP be the position vector of P in the
body-fixed system xyz whose origin is O, and whose axes x, y, and z are the
principal axes of inertia of the satellite. As shown in the preceding figure, there
results
r ¼ rO þ q
where rO is directed from G to O, that is, oppositely to the zo-axis of the orbital
system xoyozo defined in Sect. 9.1. The vectors rO and q have the following com-
ponents in the body-fixed system xyz
rO ¼ ðrO sin Uy Þux ðrO cos Uy sin Ux Þuy ðrO cos Uy cos Ux Þuz
q ¼ xux þ yuy þ zuz
where ux, uy, and uz are the unit vectors along, respectively, x, y, and z.
The force of gravity exerted by a spherical Earth of gravitational parameter lE on
the infinitesimal mass dm placed in P is
lE dm
df ¼ r
r3
9.5 Moments Due to Gravity Gradient 1145
and the moment of this force about the centre of mass O of the satellite is
lE dm
dM O ¼ q df ¼ q r
r3
Since r = rO + q, as shown in the preceding figure, then
lE dm l dm
dM O ¼ 3
q ðrO þ qÞ ¼ E 3 q rO
r r
Since the square of the magnitude of r is

" 2 #
q rO q
r ¼ ðrO þ qÞ ðrO þ qÞ ¼
2
r 2O þ q þ 2rO q ¼
2
rO2 1þ þ2 2
rO rO
then the scalar quantity 1/r3 can be expressed as follows

" 2 #32
1 1 q rO q
¼ 1þ þ2 2
r 3 r 3O rO rO
By neglecting the term (q/rO)2 in comparison with the other terms, and
expanding 1 + 2(rO q)/r2O in series, there results

1 1 rO q
1 3
r 3 r 3O r 2O
By substituting the preceding expression of 1/r3 into the expression of the

infinitesimal moment dMO, and integrating over the total mass m of the satellite,
there results
Z
lE rO q
MO ¼ 13 ðq rO Þdm
rO3 r 2O
SinceRthe origin of the system xyz is the centre of mass O of the satellite, then the
integral qdm vanishes, and therefore the preceding expression reduces to
Z
3lE rO q
MO ¼ ðq rO Þdm ¼ MOx ux þ MOy uy þ MOz uz
r 3O rO2
By executing the scalar product rO q and the vector product q rO, there
results
rO q ¼ rO ðx sin Uy y cos Uy sin Ux z cos Uy cos Ux Þ
and
q rO ¼ r O ½ðy cos Uy cos Ux þ z cos Uy sin Ux Þux

þ ðz sin Uy þ x cos Uy cos Ux Þuy þ ðx cos Uy sin Ux y sin Uy Þuz
Therefore, the components of the vector MO are
3
MOx ¼ x2o Iz Iy sinð2Ux Þ cos2 Uy
2
3
MOy ¼ x2o ðIz Ix Þ sin 2Uy cos Ux
2
3
MOz ¼ x2o Ix Iy sin 2Uy sin Ux
2
where xo = (lE/r3O)½ is the angular velocity of the satellite in a circular orbit of

radius rO. The preceding expressions show that the moment due to the gravity
gradient vanishes when Uy and Ux are equal to 0 or p/2.
In case of Ux being equal to zero, the preceding equations reduce to
MOx ¼ 0
3
MOy ¼ x2o ðIz Ix Þ sin 2Uy
2
MOz ¼ 0
and the two correspondent attitudes of equilibrium (for Ix > Iz) are shown in the
following figure.
It is important to consider the stability of the equilibrium in the two cases. The
condition of static stability is
dMOy
\0
dUy
in other words, in order for the attitude of a satellite to be statically stable, a small
increase in Uy must give rise to a restoring torque.
The derivative of the function MOy = 3/2x2o(Iz − Ix) sin(2Uy) with respect to Uy
yields
dMOy
¼ 3x2o ðIz Ix Þ cos 2Uy
dUy
For a rigid body such that Ix = Iy > Iz, the expression written above is negative
when Uy = 0, and positive when Uy = p/2. In other words, the stable attitude of a
cylindrical satellite is the attitude in which the axis of smaller moment of inertia is
directed vertically. The stable attitude (Uy = 0) is shown on the left-hand side of the
preceding figure. By contrast, the attitude (Uy = p/2) shown on the right-hand side
of the same figure is not statically stable, because in that case dMOy/dUy is positive.
The maximum moment MOy(max) due to gravity gradient corresponds to Uy = p/4,
in which case sin(2Uy) = 1. The magnitude of the gravity force acting on a satellite of
mass m is
lE m
fG ¼
rO2
The moments of inertia may also be expressed in terms of the radii of gyration,
as follows
Ix ¼ mkx2
Iz ¼ mkz2
Therefore, the maximum moment (Uy = p/4) due to gravity gradient may also be
expressed as follows
3 3l 3 k 2 kx2
MOyðmaxÞ ¼ x2o ðIz Ix Þ ¼ 3E m kz2 kx2 ¼ fG z
2 2 rO 2 rO
This moment may be considered as caused by the gravity force, fG, acting on the
satellite at a distance c from the centre of mass O of the satellite itself. In other
words, by equating fGc to the preceding expression of MOy(max), there results
3 k2 kx2
fG c ¼ fG z
2 rO
hence
3 kz2 kx2
c¼
2 rO
As an example, for the cylindrical satellite described above, Wiggins [55] has
given the following values of the moments of inertia
Ix ¼ Iy ¼ 1:28 1011 g cm2 ¼ 12800 kg m2

Iz ¼ 4:08 109 g cm2 ¼ 408 kg m2
These values make it possible to compute the maximum moment (Uy = p/4) due
to gravity gradient at the altitude of interest, by means of the expression shown
above
3 3 GmE
MOyðmaxÞ ¼ x2o ðIz Ix Þ ¼ ðIz Ix Þ
2 2 rO3
For example, at an altitude of 185 km from the surface of the Earth, there results
h . 3 i
MOyðmaxÞ ¼ 1:5 6:67384 1011 5:972 1024 6:371 106 þ 1:85 105
ð408 12800Þ ¼ 0:02629 N m ¼ 2:629 105 dyne cm
where 6.371 106 m and 5.972 1024 kg are, respectively, the mean radius and
the mass of the Earth. The value of MOy(max) is negative, because, for the satellite
considered, Iz < Ix.
The torque due to gravity gradient can be used to stabilise a satellite in a desired
attitude with respect to the Earth, as will be shown in the following sections. Suffice
it for the moment to say that a satellite revolving around the Earth can be passively
stabilised by extending vertically its long axis, for example, by means of a boom.
By so doing, the part of an orbiting structure which is nearer to the Earth is subject
to a greater gravitational attraction than that which acts on the other part. As a result
of this difference of gravitational attraction, an elongated satellite tends to align its
axis of minimum moment of inertia in the vertical direction, as has been shown for
the cylindrical satellite described above. This method of stabilisation was tested in
flight by means of the DODGE (Department Of Defence Gravity Experiment)
satellite, NSSDC/COSPAR ID: 1967-066F, shown in the following figure, which is
due to the courtesy of the United States Air Force [51].
This satellite was launched on the 1st of July 1967 from Cape Canaveral for the
primary purpose of testing biaxial and triaxial gravity-gradient stabilisation techniques
at near-synchronous altitudes (some orbital data are hP = 33270 km, hA = 33659 km,
T = 21h:58m:54s, i = 6°12′, e = 0.00488 at 1967-07-01, 12h:12m:00s UTC, according
to Ref. [32]).
Secondary objectives were measurements of the magnetic field of the Earth at
those altitudes and black-and-white and colour TV photography of the entire disc of
the Earth. This satellite carried 10 booms which were radio-commanded to extend
or retract along three different axes. It also carried a number of commandable
magnetic-damping devices and two TV cameras to determine the alignment of the
satellite. One of these cameras provided the first colour pictures of the full Earth.
The satellite was successfully stabilised 12 days after launch by means of the
gravity-gradient booms and libration dampening systems. It was oriented with its
base and mast directed towards the centre of the disc of the Earth. The mission was
successful and proved the feasibility of achieving triaxial gravity-gradient stabili-
sation at near-synchronous altitudes by using passive and semi-passive techniques.
The magnitude of the torque due to gravity gradient is proportional to 1/r3, as is

the case with the torque due to electromagnetic induction.
9.6 Moments Due to Micrometeorites
These moments are generated by micrometeorites or cosmic dust. According to

Wertz [54], such moments have been found negligible in comparison with the other
moments, except perhaps in some regions of the Solar System, such as those inside
the rings of Saturn.
9.7 Comparison of the Magnitudes of the External

Moments
A comparison concerning the magnitudes of the moments indicated in the pre-

ceding sections has been made by Gerlach [11] with reference to the cylindrical
satellite described in Sect. 9.2. The results found by him are summarised below.
They indicate that the moments due to aerodynamic forces are dominant for
satellites orbiting at altitudes less than about 500 km.
In case of a coil consisting of one winding and carrying a current of 1 A, used
for generating an electromagnetic moment, this moment has at all altitudes nearly
the same magnitude as the gravity-gradient moment generated in a satellite whose
axis of minimum moment of inertia deviates by 1° from the vertical [11]. Of course,
the magnitude of the electromagnetic moment may be increased by using coils
consisting of a greater number of windings or carrying a more intense current.
Above an altitude of 1000 km, the aerodynamic moment is smaller than the
moment due to solar radiation pressure, and above an altitude between 2000 and
3000 km the aerodynamic moment decreases to such an extent as to be comparable
with the moment due to micrometeorites.
At altitudes equal to or greater than nearly 1000 km, the moments to be taken into
account are only those due to gravity gradient, electromagnetic induction, and solar
radiation pressure, the last of which is initially smaller than the other two. At an altitude
of about 10000 km, all of the three moments cited above have nearly the same
magnitude, and at greater altitudes the moment due to solar radiation pressure becomes
dominant. Again, these results hold for the cylindrical satellite described previously.
The moment due to solar radiation pressure contains a periodic component,
which has to be taken into account for a satellite meant to remain at a fixed attitude
with respect to the Earth, unless the solar radiation comes perpendicularly to the
orbital plane of the satellite. Apart from this particular case, the periodic component
of the moment due to solar radiation pressure excites the satellite about the axis
perpendicular to the orbital plane. This fact may generate large deviations from the
desired attitude, even when the moment due to solar radiation pressure is smaller
than the stabilising moment due to gravity gradient.
9.7 Comparison of the Magnitudes of the External Moments 1151
Finally, a satellite of elongated shape is more subject than one of spherical shape
to perturbing moments [11].
9.8 Single-Spin and Dual-Spin Stabilisation of Satellites
A simple method for stabilising an artificial satellite in attitude is based on the

rotational motion of rigid bodies, which has been considered in Chap. 7. In order to
apply the theory discussed previously, a satellite is modelled as a rigid body con-
taining no moving parts.
Under this hypothesis, a satellite spinning about one of its principal axes of
inertia and not subject to external torques maintains its attitude in space.
However, as has been shown in the preceding sections of the present chapter,
there are perturbing torques which can modify the attitude of a satellite, when its
rotational torque-free motion is not stable.
The attitude of a spacecraft can be made stable by setting the vehicle spinning.
By so doing, the gyroscopic effect of the rotating mass of the vehicle is the sta-
bilising mechanism. An example of a single-spin-stabilised satellite is provided by
the Lunar Prospector (NSSDCA/COSPAR ID: 1998-001A), which is shown in the
following figure, due to the courtesy of NASA [33].
Other examples are provided by the Pioneer 10 and 11 spacecraft meant to

explore the outer Solar System, and the Galileo Jupiter orbiter spacecraft and its
atmospheric probe. In this type of stabilisation, the onboard thrusters are fired only
occasionally, in order to make desired changes in spin rate, or in the spin-stabilised
attitude. In the case of Galileo’s Jupiter atmospheric probe, and the Huygens Titan
probe, the proper attitude and spin are initially imparted by the mother ship [34].
The mathematical concepts discussed in Chap. 7 are now applied to the case of a
spinning satellite, as will be shown below. In case of a torque-free motion, the
attitude of a satellite is governed by the three Euler equations of Sect. 7.6:

Ix xx Iy Iz xy xz ¼ 0
Iy xy ðIz Ix Þxz xx ¼ 0

Iz xz Ix Iy xx xy ¼ 0
where x, y, and z are the three principal axes of inertia of the satellite.
When xy and xz are, both of them, equal to zero, then the first of these equations
shows that x`x is also equal to zero, that is, xx is constant. Likewise, the second and
the third of these equations show that xy is constant when xz = xx = 0, and xz is
constant when xx = xy = 0. In other word, a steady rotation of the satellite is
possible about each of its principal axes of inertia. As to the stability of these
rotations, let us suppose that the satellite shown in the following figure possesses, at
the initial time t = 0, a variation of angular velocity x`x0 about the x-axis.
Let us suppose that, as a result of perturbing torques, the satellite also has an
infinitesimal variation of angular velocity dx`x, dx`y, and dx`z about each of its
three principal axes of inertia, so that
xx ¼ xx0 þ dxx
xy ¼ dxy
xz ¼ dxz
The Euler equations written above, neglecting the products of infinitesimal

terms, become
Ix xx ¼ 0
Iy xy ¼ ðIz Ix Þxz xx0

Iz xz ¼ Ix Iy xy xx0
9.8 Single-Spin and Dual-Spin Stabilisation of Satellites 1153
The first of these equations indicates that xx is constant. In other words, the
angular velocity xx about the axis x of initial rotation neither increases nor
decreases, but remains constant after the perturbation.
The second of the Euler equations written above, differentiated with respect to
time, yields
Iz Ix
xy ¼ xx0 xz
Iy
The third equation, solved for x`z, yields
Ix Iy
xz ¼ xy xx0
Iz
By substituting this expression of x`z into the expression of x``y written above,
there results

Iz Ix Iy Ix
xy þ ðxx0 Þ2 xy ¼ 0
Iy Iz
xy þ k2 xy ¼ 0
where
Iz Ix Iy Ix
k2 ¼ ðxx0 Þ2
Iy Iz
xz þ k2 xz ¼ 0
The positive or negative sign of the quantity k2 depends on the positive or

negative sign of the product (Iz − Ix)(Iy − Ix).
When k2 is positive, then the two differential equations x``y + k2xy = 0 and
x``z + k2 xz = 0 written above have the following periodic solutions

xy0
xy ¼ xy0 cosðktÞ þ sinðktÞ
k

xz0
xz ¼ xz0 cosðktÞ þ sinðktÞ
k
which indicate undamped harmonic oscillations. These oscillations do not grow, but
remain bounded, and therefore the original steady motion of the satellite is stable
after the perturbation.
When k2 is negative, then the motion of the satellite after the perturbation
consists of two aperiodic parts: a damped (or decreasing) motion and an undamped
(or increasing) motion. Therefore, when k2 is negative, then the original steady
motion of the satellite is unstable after the perturbation.
Because of the definition of k2 given above, the condition for stability is satisfied
in the two following cases:
(a) Iz > Ix and Iy > Ix; or
(b) Iz < Ix and Iy < Ix.
In the first case, Ix is the smallest of the three moments of inertia; in the second
case, Ix is the greatest of them. These cases are illustrated on, respectively, the
left-hand side and the right-hand side of the following figure.
The results shown above are in agreement with those of Sect. 7.7, concerning an
axially symmetric rigid body. However, since a satellite is not in practice a rigid
body, due to the presence of elastic vibrations induced by gyroscopic action, then
dissipation of energy takes place within the satellite. Thomson [50] has shown that,
when this dissipation of energy is taken into account, the conditions for stability
discussed above must be revised, because even a small deviation of the spin axis of
a satellite may grow into a large one. This fact can eventually lead to a complete
change of the attitude of a satellite. Therefore, in case of bodies subject to dissi-
pation of energy, only a rotation about the principal axis of maximum moment of
inertia is stable. By contrast, when a rotation takes place about the principal axis of
minimum moment of inertia, then the equilibrium is unstable, as will be shown
below.
Let x be the spin axis of an axially symmetric satellite, as shown in the following
figure. In other words, we assume Iy = Iz 6¼ Ix.
As has been shown in Sect. 7.2, the magnitude of the moment-of-momentum

vector hO of this satellite is
h 2 2 i12 h i12
hO ¼ ðIx xx Þ2 þ Iy xy þ Iy xz ¼ ðIx xx Þ2 þ Iy2 x2y þ x2z
In case of no moments acting on the satellite, the moment-of-momentum vector

hO is constant in magnitude and direction.
As has been shown in Sect. 7.3, the kinetic energy possessed by a satellite due to
its rotational motion is
1 2 1h i
Trot ¼ Ix xx þ Iy x2y þ Iy x2z ¼ Ix x2x þ Iy x2y þ x2z
2 2
The two equations written above, combined together, yield

h2O 2Trot Iy ¼ Ix2 Ix Iy x2x
As has been shown in Sect. 7.7, when a satellite rotates not only about the x-axis
but also about the axes y and z, the moment-of-momentum vector hO has not the
same direction as the x-axis. With reference to the preceding figure, let h be the
angle which the moment-of-momentum vector hO forms with the x-axis. As shown
in this figure, the projection of hO onto the x-axis is
hOx ¼ Ix xx ¼ hO cos h
Because of the initial perturbation, h differs from zero. In order for the original
steady motion of the satellite to be asymptotically stable, h must return to zero after
the perturbation. In mathematical terms, the condition of stability is
h\0
The equation

h2O 2Trot Iy ¼ Ix2 Ix Iy x2x
derived above may also be rewritten as follows

Iy
h2O 2Trot Iy ¼ h2O cos2 h 1
Ix
In the preceding equation, hO is constant in time, but Trot may decrease due to
internal dissipation of energy caused, for example, by hysteresis of the deforming
material of which the satellite is made. The preceding equation, solved for Trot,
yields

1 h2O Iy
Trot ¼ 1 1 cos2 h
2 Iy Ix
which in turn, differentiated with respect to time, yields

Ix Iy
ðTrot Þ ¼ h2O hcos h sin h
Ix Iy
In case of dissipation of energy, (Trot)` can only be negative. Therefore, if Ix > Iy,
then h` is negative, and if Ix < Iy, then h` is positive. In other words, the condition
(h` < 0) to be satisfied for the stability of an axially symmetric (Iy = Iz) satellite
rotating about its x-axis is Ix > Iy. This means that the spin axis of a satellite subject
to internal dissipation of energy must be the principal axis of maximum moment of
inertia. When this condition is satisfied, then any small angle h, which may exist
initially between the moment-of-momentum vector hO and the spin axis x of the
satellite, decreases to zero. By contrast, when the spin axis of the satellite is the
principal axis of minimum moment of inertia, then the angle h between the
moment-of-momentum vector hO and the spin axis x of the satellite increases after
any small perturbation. Due to the increase of this angle, the satellite rotates no
more about its initial spin axis, but about an axis perpendicular to it. Such was the
case with the Explorer I (NSSDC/COSPAR ID: 1958-001A) satellite, launched in
1958, which changed its initial axis of rotation. The Explorer I satellite is shown in
the following figure, due to the courtesy of NASA [35].
This happened because the prolate body of the satellite had been designed to
spin about its principal axis of minimum moment of inertia. Instead of doing so, the
satellite changed its previous attitude into a state of minimum kinetic energy, that is,
into a state of a flat spin about the transverse axis. This was deduced from the
modulation of the received signal, which produced periodic fade-outs of the signal
[36]. In the case of the Explorer I, the energy dissipation was due to the flexible
antennas of the satellite shown above, which dissipated energy through hysteresis in
the material [11].
A satellite stabilised by means of the simple technique described above must be
an oblate (that is, fat and short) body, spinning about its axis of maximum moment
of inertia. However, the diameter of a spacecraft is limited by the cross section of
the upper stage of the vehicle used to launch it, and its length is limited by the
requirement of stability. Therefore, an oblate spinning satellite cannot take
advantage to the full extent of the volume made available by launch vehicles, which
are prolate for aerodynamic reasons. In addition, if the communication antennas of
a satellite were fixed to and spinning together with the main body, the communi-
cation efficiency of this satellite would be very low. This is because, in that case, the
antenna beam would scan the disc of the Earth for a short time in each spin period,
and be directed elsewhere for the remaining time, when no use could be made of the
emitted energy. In order to avoid this, the antenna is mechanically or electronically
de-spun.
A dual-spin-stabilised satellite, also known as a gyrostat, makes it possible to
avoid the undesirable effects indicated above. It consists of an axially symmetric
rotor and a smaller axially symmetric platform, aligned with the rotor along a
common longitudinal spin axis at a bearing, as shown in the following figure, due to
the courtesy of CPinterSEP [7], which shows the Boeing HS 376 satellite. This
geostationary satellite, launched in 1980, was spin-stabilised by means of its rotor,
which spun at a rate of about 50 rotations per minute about its symmetry axis,
whereas its platform spun at a rate of 1 rotation per orbit, in order for the antenna to
point constantly at a desired place of the Earth.
With reference to the scheme shown in the following figure, let O, P, and R be
the centres of mass of, respectively, the whole satellite, the platform, and the rotor.
In case of a satellite made of a homogeneous material, these points are placed along
its axis of symmetry (x). This axis is also one of the three principal axes of inertia x,
y, and z of the satellite.
Since the satellite is axially symmetric, then its moments of inertia Iy and Iz about
the other two axes y and z are equal (Ix 6¼ Iy = Iz). The platform and the rotor have
their own components of angular velocity, respectively, xp and xr, along the
longitudinal spin axis x, which is directed along the unit vector ux.
The angular spin of the platform is much lower than that of the rotor. The whole
satellite can be modelled as a rigid body, as far as the transverse angular velocity
(xt) is concerned. In other words, the platform and the rotor have the transverse
angular velocity vector xt in common. This vector lies in the yz plane, which passes
through O and is perpendicular to the spin axis x.
An electric motor is integrated into the axle bearing, in order to counteract the
frictional torque, which would otherwise reduce to zero the difference of angular
velocity between the rotor and the platform. In that event, the spinning satellite
would become unstable, due to its shape, which is usually prolate.
In case of a communication satellite revolving about the Earth in a geostationary
orbit, the platform spins at a rate of one cycle per sidereal day, in order for its
antennas to point towards the Earth. The rotor spins at a rate of about one cycle per
second. The spin axis of the satellite is perpendicular to the orbital plane.
In case of sensor and antenna systems which require inertial pointing, the
platform is completely de-spun. Such was the case with the Galileo Jupiter probe,
the first dual-spin-stabilised spacecraft used in an outer-planet mission. The rotor of
the Galileo spacecraft spun at a rate of three cycles per minute.
The total moment of momentum of a dual-spin-stabilised satellite about the
centre of mass O is the vector sum of the moment-of-momentum vector relative to
the rotor (subscript r) and the moment-of-momentum vector relative to the platform
(subscript p), as follows
hO ¼ hOr þ hOp
The moment-of-momentum vector relative to the rotor about the centre of mass
O of the satellite is
hOr ¼ Ixr xr ux þ Itr xt
where Ixr is the moment of inertia of the rotor about the x-axis, and Itr = Iyr = Izr is
the moment of inertia of the rotor about the transverse axis.
Likewise, the moment-of-momentum vector relative to the platform about the
centre of mass O of the satellite is
hOp ¼ Ixp xp ux þ Itp xt
where Ixp is the moment of inertia of the platform about the x-axis, and
Itp = Iyp = Izp is the moment of inertia of the platform about the transverse axis. The
total moment-of-momentum vector of the satellite about the centre of mass O is

hO ¼ hOr þ hOp ¼ ðIxr xr þ Ixp xp Þux þ Itr þ Itp xt ¼ ðIxr xr þ Ixp xp Þux þ It xt
where It = Itr + Itp is the moment of inertia of the whole satellite, including the rotor
and the platform, about the transverse axis.
Since the two components of the vector hO indicated in the preceding equation
are taken along two perpendicular axes, then the square of the magnitude of this
vector is
h2O ¼ ðIxr xr þ Ixp xp Þ2 þ It2 x2t
As has been shown in Sect. 7.7, in case of torque-free motion of an axially

symmetric rigid body, hO is constant in time, and therefore the first derivative of h2O
with respect to time is equal to zero:
2ðIxr xr þ Ixp xp ÞðIxr xr þ Ixp xp Þ þ It2 ðx2t Þ ¼ 0
The preceding equation, solved for (x2t )`, yields

Ixr xr þ Ixp xp
x2t ¼ 2 Ixr xr þ Ixp xp
It2
The rotational kinetic energy of the whole satellite is the sum of the rotational
kinetic energy of the rotor (superscript r) and the rotational kinetic energy of the
platform (superscript p), as follows
ðrÞ ðpÞ
Trot ¼ Trot þ Trot ¼ 1=2Ixr x2r þ 1=2Ixp x2r þ 1=2It x2t
The preceding expression, differentiated with respect to time, yields
ðTrot Þ ¼ Ixr xr xr þ Ixp xp xp þ 1=2It ðx2t Þ
which in turn, solved for (x2t )`, yields
2 ðTrot Þ Ixr xr xr Ixp xp xp

xt ¼ 2
It
The derivative, (Trot)`, of the rotational kinetic energy with respect to time is the
sum of the power, P(r), dissipated in the rotor and the power, P(p), dissipated in the
platform, that is, (Trot)` = P(r) + P(p), and consequently
PðrÞ þ PðpÞ Ixr xr xr Ixp xp xp

ðx2t Þ ¼ 2
It
By comparing this expression of (x2t )` with the other expression of (x2t )` written
above, that is, with

Ixr xr þ Ixp xp
x2t ¼ 2 Ixr xr þ Ixp xp
It2
there results

PðrÞ þ PðpÞ Ixr xr xr Ixp xp xp Ixr xr þ Ixp xp
2 ¼ 2 Ixr xr þ Ixp xp
It It2
The preceding equation, solved for P(r) + P(p), yields
Ixr
Ixp

PðrÞ þ PðpÞ ¼ ðIt Ixr Þxr Ixp xp xr þ It Ixp xp Ixr xr xp
It It
Likins [25] and Curtis [8] identify the first term on the right-hand side with the
power dissipated in the rotor, and the second term on the right-hand side with the
power dissipated in the platform. In other words, they set
Ixr
Ixp

PðrÞ ¼ ðIt Ixr Þxr Ixp xp xr PðpÞ ¼ It Ixp xp Ixr xr xp
It It
The two equations written above, solved for, respectively, x`r and x`p, yield
It PðrÞ It PðpÞ
xr ¼
xp ¼

Ixr ðIt Ixr Þxr Ixp xp Ixp It Ixp xp Ixr xr
By substituting these expressions into the preceding equation
PðrÞ þ PðpÞ Ixr xr xr Ixp xp xp

x2t ¼ 2
It
there results
" #
2 PðrÞ PðpÞ xp
x2t ¼ þ Ixr þ Ixp
It Ixp xxpr ðIt Ixr Þ Ixr It Ixp xxpr xr
As has been shown above, in case of a dual-spin geostationary satellite used for
communications, the value of the ratio xp =xr is such that
xp 2p radians=sidereal day
¼ ¼ 0:000011606
xr 2p rad=s
In case of an interplanetary dual-spin interplanetary spacecraft, the value of the

ratio xp/xr is exactly equal to zero. Hence, for a de-spun platform such that xp/xr is
nearly or exactly equal to zero, the preceding equation reduces to

2 2 Ixr ðr Þ ðpÞ
xt ¼ P þP
It Ixr It
When the rotor is an oblate body (Ixr > It), then (x2t )` < 0, because P(r) and P(p)
are both of them less than zero, and therefore the attitude of the whole satellite is
unconditionally stable.
On the other hand, in many practical cases, the rotor is a prolate body (Ixr < It),
and therefore
Ixr PðrÞ
[0
Ixr It
In these cases, (x2t )` is less than zero (and therefore the attitude of the whole
satellite is stable) only when

Ixr PðrÞ ðpÞ

I I \ P
xr t
that is, only when the power dissipated in the platform is much greater than the
power dissipated in the rotor. In order to increase the value of |P(p)|, it is possible to
use nutation dampers, as will be shown in Sect. 9.9.
When xp/xr is nearly or exactly equal to zero, the equations written above
It PðrÞ It PðpÞ
xr ¼
xp ¼

Ixr ðIt Ixr Þxr Ixp xp Ixp It Ixp xp Ixr xr
It PðrÞ It PðrÞ
xr ¼
xp Ixr xr ðIt Ixr Þ
Ixr xr ðIt Ixr Þ Ixp
xr
It PðpÞ I PðpÞ
xp ¼ t
xp Ixp Ixr xr
Ixp xr It Ixp Ixr
xr
which shows that x`r and x`p have opposite signs. Hence, if xr > 0, then the spin
velocity of the rotor decreases and the spin velocity of the platform increases,
because P(r) and P(p) are both of them less than zero. This fact shows once again the
necessity of a motor on the shaft connecting the platform with the rotor, in order to
prevent the attitude of the whole satellite from becoming unstable in case of
xp = xr.
9.9 Nutation Dampers
These devices are passive means for dissipating kinetic energy associated with
undesired nutations of a single-spin- or dual-spin-stabilised satellite. A common
type of such devices is a tube filled with a viscous fluid (e.g. silicone oil) and a mass
attached to springs, as shown in the following figure.
A nutation damper may also contain only fluid which occupies part of the
volume made available by the tube, so as to slosh around. The kinetic energy
associated with the motion of nutation is dissipated by friction in the viscous fluid.
With reference to the following figure, let O and G be, respectively, the centre of
mass of the satellite alone, considered as a rigid body, and the centre of mass of the
system comprising the satellite and the nutation damper.
The point G is placed along the straight line joining the point O with the point
occupied by the particle of mass m contained in the tube. The direction of this tube
is parallel to the axis (z) of symmetry of the satellite, so that the mass m can only
move in the direction of the z-axis, and lies in the xy plane when the springs are not
deformed. The reference system xyz has its origin in the centre of mass O of the
satellite and is attached to the body. The axes x, y, and z are the principal axes of
inertia of the body.
Following Curtis [8], we want to write the equation of motion for the particle of
mass m in the system xyz, which is not an inertial system. For this purpose, it is
necessary to consider the external accelerations acting upon the particle, the relative
acceleration (arel) which the particle possesses in the reference system xyz, and the
fictitious accelerations, as will be shown below.
Let R be the radius of the cylindrical satellite. The vector r, which indicates the
position of the particle of mass m in the reference system xyz attached to the
satellite, has the following components
r ¼ Rux þ 0uy þ zm uz
where zm is a function of time, and ux, uy, and uz are the unit vectors along the axes,
respectively, x, y, and z of the body-fixed system. Therefore, the velocity and
acceleration vectors of the particle, in the reference system xyz, are respectively
vrel ¼ r ¼ zm uz
arel ¼ r ¼ zm uz
Let
x ¼ xx ux þ xy uy þ xz uz
be the angular velocity vector of the satellite in an inertial reference system. The
derivative of this vector with respect to time in the body-fixed system xyz is
The acceleration vector of the particle of mass m in the inertial system is
a ¼ aO þ arel þ 2x vrel þ x ðx rÞ þ x r
where 2x vrel = 2x r` is the Coriolis acceleration, x (x r) is the cen-

trifugal acceleration, and x` r is the Euler acceleration, all of which are fictitious
accelerations. In addition, aO is the acceleration vector acting on O in the inertial
reference system. The components of the vectors arel, vrel, x, r, and x` along the
axes x, y, and z are specified above. By executing the operations indicated on the
right-hand side of the preceding equation, we obtain
9.9 Nutation Dampers 1165
x r ¼ zm xy ux þ ðRxz zm xx Þuy Rxy uz

x ðx rÞ ¼ ðRx2y Rx2z þ zm xx xz Þux þ ðRxx xy þ zm xy xz Þuy
þ ðRxx xz zm x2x zm x2y Þuz
2x vrel ¼ 2zm xy ux 2zm xx uy þ 0uz
a ¼ ðaOx þ 2zm xy Rx2y Rx2z þ zm xx xz þ zm xy Þux
þ ðaOy 2zm xx þ Rxx xy þ zm xy xz þ Rxz zm xx Þuy
þ ðaOz þ zm þ Rxx xz zm x2x zm x2y Rxy Þuz
Therefore, the components of the acceleration vector a acting on the particle of

mass m along the axes x, y, and z may be written as follows
ax ¼ aOx þ 2zm xy Rðx2y þ x2z Þ þ zm xx xz þ zm xy

ay ¼ aOy 2zm xx þ Rxx xy þ zm xy xz þ Rxz zm xx
az ¼ aOz þ zm þ Rxx xz zm ðx2x þ x2y Þ Rxy
Let us consider now the forces acting on the particle of mass m. With reference
to the following figure, let wx, wy, and wz be the three components of the force of
gravity w along the axes, respectively, x, y, and z. Let nx and ny be the components
of the force of contact n between the particle and the walls of the tube, supposed to
be smooth. Let kzm and cz`m be, respectively, the force exerted by the spring and the
viscous drag exerted by the fluid contained in the tube. The coefficients k and c are,
respectively, the elastic constant of the spring and the damping constant of the fluid.
Therefore, the three components of the total force f acting on the particle are
fx ¼ wx nx
fy ¼ wy ny
fz ¼ wx kzm czm
The directions of these components are chosen arbitrarily.

By using the preceding equations, which express the components of, respec-
tively, the acceleration vector and the force vector acting on the particle of mass m,
the second principle of dynamics (f = ma) may be written as follows
fx ¼ wx nx ¼ maOx þ 2mzm xy mRðx2y þ x2z Þ þ mzm xx xz þ mzm xy

fy ¼ wy ny ¼ maOy 2mzm xx þ mRxx xy þ mzm xy xz þ mRxz mzm xx
fz ¼ wx kzm czm ¼ maOz þ mzm þ mRxx xz mzm ðx2x þ x2y Þ mRxy
These equations of motion may be expressed in the following form
nx ¼ mRðx2y þ x2z Þ mzm xy mzm xx xz 2mzm xy þ fwx maOx g

ny ¼ mRxz mRxx xy þ mzm xx mzm xy xz þ 2mzm xx þ wy maOy
mzm þ czm þ ½k mðx2x þ x2y Þzm ¼ mRðxy xx xz Þ þ fwz maOz g
The three terms in curly brackets are equal to zero when the acceleration of
gravity acting on the particle of mass m is the same as that acting on the centre of
mass O of the satellite. This holds when the gravity gradient between these points is
negligible. In addition, when the mass m contained in the nutation damper is much
smaller than the mass of the satellite considered as a rigid body, then the effect of
m on the rotary motion of the satellite is negligible.
When the external torques acting on the axially symmetric satellite are equal to
zero, then the components of the angular velocity vector x along the principal axes
of inertia x, y, and z of the satellite are those of Sect. 7.7:
xx ¼ x sin a cosðktÞ
xy ¼ x sin a sinðktÞ
xz ¼ x cos a
where a is the constant angle which the angular velocity vector x forms with the
principal axis z, x = (x2x + x2y + x2z )½, and k = (Iz/Ix − 1)xz = (Iz/Ix − 1)x cos a.
Of course, the presence of either the sine or the cosine of kt in the functions of
time which express xx and xy depends on the initial conditions (t = 0).
The derivatives of xx, xy, and xz with respect to time are
xx ¼ xk sin a sinðktÞ

xy ¼ xk sin a cosðktÞ
xz ¼ 0
These functions of time and their derivatives, introduced into the three scalar
equations of motion written above, yield
nx ¼ mR½x2 sin2 a sin2 ðktÞ þ x2 cos2 a mzm xðk þ x cos aÞ sin a cosðktÞ
2mzm x sin a sinðktÞ
ny ¼ mRx2 sin2 a sinðktÞ cosðktÞ mzm xðk þ x cos aÞ sin a sinðktÞ
þ 2mzm x sin a cosðktÞ
mzm þ czm þ ðk mx2 sin2 aÞzm ¼ mRxðk x cos aÞ sin a cosðktÞ
The third of the equations written above is the differential equation which
governs the motion of a damped harmonic oscillator, having one degree of freedom
and driven by a sinusoidal force.
In other words, the motion of precession of the angular velocity vector x about
the principal axis z of inertia of the satellite produces a periodic force, whose
amplitude is mRx (k – x cos a) sin a, and whose angular frequency is k. This force
causes the particle of mass m to oscillate back and forth in the tube of the nutation
damper.
In order to determine the unknown function of time zm which satisfies the third
equation, we express this function in the following form
zm ¼ A cosðkt uÞ
where A is the amplitude and u is the phase lag of the oscillation of the particle with
respect to the phase of the forcing function. By differentiating the preceding
function with respect to time, there results
zm ¼ Ak sinðkt uÞ

zm ¼ Ak2 cosðkt uÞ
The function zm specified above and its derivatives z`m and z``m are introduced
into the third differential equation. This yields
Amk2 cosðkt uÞ Ack sinðkt uÞ þ Aðk mx2 sin2 aÞ cosðkt uÞ

¼ mRxðk x cos aÞ sin a cosðktÞ
By using the trigonometric identities
cosða bÞ ¼ cos a cos b þ sin a sin b

sinða bÞ ¼ sin a cos b cos a sin b
and grouping the coefficients of cos(kt) and sin(kt), there results
fA½k mðk2 þ x2 sin2 aÞ cos u þ Ack sin u mRxðk x cos aÞ sin ag cosðktÞ
þ fA½k mðk2 þ x2 sin2 aÞ sin u Ack cos ug sinðktÞ ¼ 0
In order for the preceding equality to hold at any time, the two quantities within
curly brackets must be equal to zero. This condition gives rise to the following two
algebraic equations
A½k mðk2 þ x2 sin2 aÞ cos u þ Ack sin u ¼ mRxðk x cos aÞ sin a
A½k mðk2 þ x2 sin2 aÞ sin u ¼ Ack cos u
The second of these equations yields

" #
ck
u ¼ arctan
k m k2 þ x2 sin2 a
Since tan u = ck/[k – m(k2 + x2 sin2a)], then we use the following trigono-
metric identities
tan u 1
sin u ¼ 1 cos u ¼ 1
ð1 þ tan2 uÞ 2
ð1 þ tan2 uÞ2
and compute A by substituting sin u and cos u, expressed as a function of tan u,

into the first equation. This yields
mRxðk x cos aÞ sin a

A¼n o12

2
k m k2 þ x2 sin2 a þ c2 k2
The phase lag u shows that zm is a periodic function of time, which has two
components. One of them is in phase, and the other is in quadrature (p/2 radians out
of phase) with respect to the forcing function. The oscillation of the particle
zm = A cos(kt − u) results from the superposition of these two components.
So much for the third of the three scalar differential equations written above.
The other two equations show that the components nx and ny of the force n of
contact of the particle with the tube depend only on the amplitude and angular
frequency of the motion of precession. In the absence of precession, in which case
a = 0 and therefore the angular velocity vector x is aligned with the z-axis of the
satellite, the three differential equations written above reduce to
nx ¼ mRx2
ny ¼ 0
zm ¼ 0
In this ideal case, the motion of the satellite about its centre of mass reduces to a
spin about the z-axis. In other words, there is no necessity of a nutation damper
which dissipates kinetic energy.
In practice, due to the presence of precession, a nutation damping device is

necessary to stabilise the attitude of a spinning satellite. In order to describe
mathematically this case with a better degree of approximation than that used
previously, we take account of the fact that the mass m of the particle contained in
the tube of the damper is not negligible in comparison with the mass mS of the
satellite alone. Therefore, the centre of mass O of the satellite alone does not
coincide with the centre of mass G of the system comprising the satellite and the
particle in the damper. Since the whole system cannot be considered as a rigid
body, due to the motion of the particle in the damper, we cannot use the Euler
equations, which are based on this hypothesis.
Following Curtis [8], we write the fundamental equation of the rotational motion
of the system comprising the satellite and the damper as follows
hO þ rGO ðmS þ mÞaOG ¼ M G
In the preceding equation, account is taken of the motion of the point G, in

which the total mass mS + m of the system is concentrated, with respect to the point
O. In this motion, the total mass mS + m is accelerated in the direction of the vector
aOG. The force associated with this acceleration is
f ¼ ðmS þ mÞaOG
The moment with respect to O of this force is
rGO ðmS þ mÞaOG
Therefore, the net moment MG results from the equation written above.
The moment of momentum with respect to O of the system comprising the mass
mS of the satellite and the mass m of the particle is
hO ¼ ½Ix xx ux þ Iy xy uy þ Iz xz uz þ ½r mr0
where the first of the two expressions in square brackets on the right-hand side is
the moment of momentum due to the spacecraft considered as a rigid body, and the
second expression in square brackets is the moment of momentum due to the
particle, whose velocity vector is r′. Remembering the equation of Sect. 7.1
r0 ¼ rþ x r
and the equations shown above in the present section
r ¼ vrel ¼ zm uz
there results

r0 ¼ zm uz þ ðxx ux þ xy uy þ xz uz Þ Rux þ 0uy þ zm uz
By executing the operations indicated above, we find
r0 ¼ zm xy ux þ ðRxz zm xx Þuy þ ðzm Rxy Þuz
r mr0 ¼ ½Rux þ 0uy þ zm uz m½zm xy ux þ ðRxz zm xx Þuy þ ðzm Rxy Þuz

¼ mzm ðzm xx Rxz Þux þ mðz2m xy Rzm þ R2 xy Þuy þ mRðRxz zm xx Þuz
By introducing the preceding expression into
hO ¼ ½Ix xx ux þ Iy xy uy þ Iz xz uz þ ½r mr0
there results
hO ¼ ½ðIx þ mz2m Þxx mRzm xz ux þ ½ðIy þ mR2 þ mz2m Þxy mRzm uy
þ ½ðIz þ mR2 Þxz mRzm xx uz
We compute h′O by using again the equation
h0O ¼ hO þ x hO
This yields
h0 O ¼ ½2mzm zm xx þ ðIx þ mz2m Þxx mRzm xz mRzm xz ux

þ ½2mzm zm xy þ ðIy þ mR2 þ mz2m Þxy mRzm uy
þ ½ðIz þ mR2 Þxz mRzm xx mRzm xx uz
x hO ¼ ½ðIz þ mR2 Þxy xz mRzm xx xy ðIy þ mR2 þ mz2m Þxy xz þ mRzm xz ux
þ ½ðIx þ mz2m Þxx xz mRzm x2z ðIz þ mR2 Þxx xz þ mRzm x2x uy
þ ½ðIy þ mR2 þ mz2m Þxx xy mRzm xx ðIx þ mz2m Þxx xy þ mRzm xy xz uz
Therefore, the first term of the equation h′O + rGO (mS + m) aOG = MG is
h0 O ¼ hO þ x hO ¼ hOx ux þ hOy uy þ hOz uz
where the components of the vector h′O along the axes x, y, and z are
h0 Ox ¼ ðIx þ mz2m Þxx þ ðIz Iy mz2m Þxy xz þ 2mzm zm xx mRzm xz mRzm xx xy

h0 Oy ¼ ðIy þ mR2 þ mz2m Þxy þ 2mzm zm xy mRzm þ Ix Iz þ mz2m mR2 xx xz
þ mRzm ðx2x x2z Þ
h0 Oz ¼ ðIz þ mR2 Þxz 2mRzm xx mRzm xx þ ðIy Ix þ mR2 Þxx xy þ mRzm xy xz
Let us consider the second term, that is, rGO (mS + m) aOG, of the same
equation.
The vector rGO indicates the centre of mass G of the whole system (satellite and
particle) with respect to the centre of mass O of the satellite only. This vector, by its
definition, can be expressed as follows
ðmS þ mÞrGO ¼ mS 0 þ mr
where r is the position vector of the particle of mass m with respect to the centre of
mass O of the satellite only.
Remembering that the vector r has the following components
there results
ðmS þ mÞrGO ¼ mðRux þ 0uy þ zm uz Þ

m
rGO ¼ Rux þ 0uy þ zm uz ¼ lr ¼ l Rux þ 0uy þ zm uz
mS þ m
where l indicates the ratio of the mass m of the particle to the total mass.
Therefore, the second term rGO (mS + m) aOG of the equation
h0 O þ rGO ðmS þ mÞaOG ¼ M G
rGO ðmS þ mÞaOG ¼ ½m=ðmS þ mÞr ðmS þ mÞaOG

¼ ½m=ðmS þ mÞðmS þ mÞr aOG
¼ mr aOG ¼ r maOG
The vector aOG is the acceleration of the point O with respect to the point G.
This vector can be computed as follows
aOG ¼ aGO ¼ lr

r0 ¼ rþ x r
r ¼ ðrþ x rÞþ x ðrþ x rÞ ¼ rþ x r þ x rþ x r
þ x ðx rÞ ¼ rþ x r þ 2x rþ x ðx rÞ
aOG ¼ l½rþ x r þ 2x rþ x ðx rÞ
where

r ¼ Rux þ 0uy þ zm uz ¼ 0ux þ 0uy þ zm uz

r ¼ Rux þ 0uy þ zm uz ¼ 0ux þ 0uy þ zm uz
The expression
r ¼ 0ux þ 0uy þ zm uz
is substituted below into the preceding equation, remembering that

x ðx rÞ ¼ ðRx2y Rx2z þ zm xx xz Þux þ ðRxx xy þ zm xy xz Þuy
þ ðRxx xz zm x2x zm x2y Þuz
2x r ¼ 2zm xy ux 2zm xx uy þ 0uz
This yields the components along x, y, and z of the acceleration vector of the
point O with respect to the point G, as follows
aOG ¼ ½lzm xy 2lzm xy þ lRðx2y þ x2z Þ lzm xx xz ux

þ ½lRxz þ lzm xx þ 2lzm xx lRxx xy lzm xy xz uy
þ ½lzm þ lRxy lRxx xz þ lzm ðx2x þ x2y Þuz
By using these components, we can evaluate the expression

rGO (mS + m) aOG, where rGO = lRux + 0uy + lzmuz. This expression is
rGO ðmS þ mÞaOG ¼ lm½zm Rðxz þ xx xy Þ z2m xx 2zm zm xx þ z2m xy xz ux
þ lm½ðz2m þ R2 Þxy 2zm zm xy þ zm Rðx2z þ x2x Þ
þ ðR2 z2m Þxx xz þ Rzm uy þ lm½R2 xz þ zm Rxx
þ 2Rzm xx R2 xx xy zm Rxy xz uz
The components of the vector rGO (mS + m) aOG along the axes x, y, and z are
added to the corresponding components of the vector h′O, which have been com-
puted previously. The result of each of these three additions is equal to zero,
because the three components MGx, MGy, and MGz of the net moment MG are also
equal to zero. By so doing, we obtain the following threes scalar equations

Ix xx þ ðIz Iy Þxy xz þ mð1 lÞ½z2m ðxx xy xz Þ Rzm xz þ xx xy
þ 2zm zm xx ¼ 0
½ðIy þ mR2 Þ lmR2 xy þ ½ðIx þ lmR2 Þ ðIz þ mR2 Þxx xz
þ mð1 lÞ½z2m ðxy þ xx xz Þ þ 2zm zm xy Rzm þ Rzm ðx2x x2z Þ ¼ 0

Iz þ mR2 lmR2 xz þ Iy þ mR2 ðIx þ lmR2 Þxx xy
mRð1 lÞ½zm ðxx xy xz Þ þ 2zm xx ¼ 0
A further equation concerns the motion of the particle of mass m within the tube
of the nutation damper along the z-axis. Remembering that wz, kzm, and cz`m are,
respectively, the component along z of the force of gravity, the force exerted by the
spring, and the viscous drag due to the fluid, the fourth equation is
wz kzm czm ¼ maz
where the acceleration az is expressed by the equation derived above
az ¼ aOz þ zm þ Rxx xz zm ðx2x þ x2y Þ Rxy
In this equation, the term aOz may be expressed as follows
aOz ¼ aOz aGz þ aGz ¼ aOGz þ aGz
where aOGz = −lz``m + lRx`y − lRxxxz + lzm(x2x + x2y ), and aGz = wz/m.
Therefore, the fourth equation wz − kzm − cz`m = maz can be written as follows
mð1 lÞzm þ czm þ ½k mð1 lÞðx2x þ x2y Þzm ¼ mRð1 lÞðxy xx xz Þ
We have obtained the following four differential equations


Ix xx þ ðIz Iy Þxy xz þ mð1 lÞ½z2m ðxx xy xz Þ Rzm xz þ xx xy
þ 2zm zm xx ¼ 0
½ðIy þ mR2 Þ lmR2 xy þ ½ðIx þ lmR2 Þ ðIz þ mR2 Þxx xz
þ mð1 lÞ½z2m ðxy þ xx xz Þ þ 2zm zm xy Rzm þ Rzm ðx2x x2z Þ ¼ 0

½ Iz þ mR2 lmR2 xz þ ½ Iy þ mR2 ðIx þ lmR2 Þxx xy
mRð1 lÞ½zm ðxx xy xz Þ þ 2zm xx ¼ 0
mð1 lÞzm þ czm þ ½k mð1 lÞðx2x þ x2y Þzm ¼ mRð1 lÞðxy xx xz Þ
where the unknowns are the functions of time xx, xy, xz, and zm. These nonlinear
equations are to be solved numerically, by using, for example, one of the methods
described in Chap. 6. Instead of solving these equations, we want to ascertain
whether they have, or have not, a stable solution near some point of equilibrium, as
For a better understanding of the method for searching a stable solution, some
fundamental concepts of the classical theory of stability due to Lyapunov are given
here. Following Murray et al. [29], let us consider a dynamical system governed by
the following differential equation
x0 ¼ f ðx; tÞ
with its associated initial condition x = x0 when t = t0, where x is a column vector
of n real values. The integrand function f(x, t) is supposed to satisfy the standard
conditions for existence and uniqueness of solution. A point x*, where x* is also a
column vector of n real values, is said to be a point of equilibrium for the dynamical
system defined above if f(x*, t) = 0 at any time t.
Roughly speaking, a point of equilibrium x* is locally stable, if all solutions x of
the differential equation x′ = f(x, t) which start near x* (meaning that the initial
value x0 is in some neighbourhood of x*) remain near x* at any time t > t0.
A point of equilibrium x* is said to be locally asymptotically stable, if x* is
locally stable, and in addition, all solutions x which start near x* tend to x* as
t tends to infinity. For example, from an intuitive point of view, a pendulum has a
point of equilibrium which is locally stable when the pendulum is hanging straight
down, and another point of equilibrium which is unstable when the pendulum is
placed straight up. If the motion of a pendulum is subject to damping, then its point
of locally stable equilibrium is also locally asymptotically stable.
In mathematical terms, let x = u(t, x0) be the existing and unique solution of the
differential equation x′ = f(x, t), which solution takes the value x0 when t = t0.
A point of equilibrium x* is said to be stable in the sense of Lyapunov if, for every
positive number e, there is a positive number d such that ||u(t, x0) − x*|| < e at any
time t > t0 for all x0 such that ||x0 − x*|| < d.
A point of equilibrium x* which does not satisfy the condition for stability
indicated above is said to be unstable. In this case, the natural response of a
dynamical system grows without bond when t > t0.
If d can be chosen not only so that the solution u(t, x0) is stable in the sense of
Lyapunov but also so that u(t, x0) tends to x* as t tends to infinity, then a point of
equilibrium x* is said to be asymptotically stable.
The definitions given above are local, because they describe the behaviour of a
dynamical system near a point of equilibrium. A point of equilibrium x* is said to
be globally stable if it is stable for all real initial conditions x0.
In addition to the states of stable equilibrium and unstable equilibrium, there is
another state, which is called neutrally or marginally stable equilibrium and is
defined as follows.
When a dynamical system, as a result of a small perturbation which moves it
from a point of equilibrium, has a natural response which neither decays nor grows
but remains constant or has bounded oscillations when t > t0, then that system is in
a state of neutrally stable equilibrium.
Now, let us come back to the problem concerning the stability of the solution of
the four nonlinear differential equations near some point of equilibrium. For this
purpose, we linearise these equations, assuming the satellite to spin at t = t0 with a
constant angular velocity x0 about the z-axis, and the mass m of the particle in the
tube of the damper to be at rest in the point of co-ordinates R, 0, and zm at the initial
time t = t0. In other words, let xx = 0, xy = 0, xz = x0 = constant, and zm = 0 be
the initial values taken by the four functions of time xx, xy, xz, and zm in the
absence of perturbations. In the presence of a perturbation, these values become,
respectively, xx, xy, x0 + xz, and zm.
By substituting the new values into the four equations and retaining only the
terms containing powers which are at most of the first degree, there results
Ix Dxx þ ðIz Iy Þx0 xy ¼ 0

½Ix Iz mR2 ð1 lÞx0 xx þ ½Iy þ mR2 ð1 lÞDxy mRð1 lÞðD2 þ x20 Þzm ¼ 0
½Iz þ mR2 ð1 lÞDxz ¼ 0
mRð1 lÞx0 xx mRð1 lÞDxy þ ½mð1 lÞD2 þ cD þ kzm ¼ 0
where the differential operator Dn() dn()/dtn has been used above to denote the nth
derivative of its argument in the body-fixed system. Consequently, Dxx stands for
x`x, Dxy stands for x`y, Dxz stands for x`z, Dzm stands for z`m, and D2zm stands for
z``m. The third of these equations yields Dxz = 0, that is, xz = constant.
The remaining three equations may be rewritten in matrix form as follows
Ax ¼ 0
where A is a 3 3 matrix, and x and 0 are two 3 1 column vectors. These

vectors have the following components
2 3 2 3
xx 0
x ¼ 4 xy 5 0 ¼ 405
zm 0
The matrix A has the following components

2 3
Ix D
Iz Iy x0 0
A ¼ 4 ½Ix Iz mR2 ð1 lÞx0 Iy þ mR2 ð1 lÞ D mRð1 lÞ D2 þ x20 5
mRð1 lÞx0 mRð1 lÞD mð1 lÞD2 þ cD þ k
In order for the solution of the system of linearised differential equations Ax = 0

to be stable, the determinant of the matrix A must be equal to zero. The condition
det(A) = 0 gives rise to the following equation
a4 D4 þ a3 D3 þ a2 D2 þ a1 D þ a0 ¼ 0
whose coefficients are expressed below
a4 ¼ mð1 lÞIx Iy
a3 ¼ cIx ½Iy þ mR2 ð1 lÞ

a2 ¼ k½Iy þ mR2 ð1 lÞIx þ mð1 lÞ½ðIx Iz Þ Iy Iz mR2 ð1 lÞIx x20

a1 ¼ cf½Ix Iz mR2 ð1 lÞ Iy Iz gx20

a0 ¼ kf½Ix Iz mR2 ð1 lÞ Iy Iz gx20 þ ½ Iy Iz ð1 lÞ2 m2 R2 x40
According to the Routh–Hurwitz criteria, the motion represented by the system

of three linearised equations in matrix form Ax = 0 (and hence by the fourth-degree
equation a4D4 + a3D3 + a2D2 + a1D + a0 = 0) is asymptotically stable if and only
if the following conditions are satisfied:
ai [ 0 ði ¼ 0; . . .; 4Þ
a3 a 2 [ a4 a 1
a3 a2 a1 [ a4 a21 þ a23 a0
As an example of application, let us consider the stability of a cylindrical satellite

spinning about its axis (z) of maximum moment of inertia at an angular velocity
x0 = 1 cycle (2p radians) per second. Let Ix = Iy = 300 kg m2 and Iz = 500 kg m2
be the moments of inertia about the principal axes of inertia x, y, and z of the
satellite. Let R = 1 m be the radius of the satellite, m = 10 kg be the mass of the
particle contained in the nutation damper, l = 0.01 be the ratio of the mass m of the
particle to the total mass m + mS, k = 10000 N/m be the constant of the spring, and
c = 150 Ns/m be the coefficient of damping.
Since the satellite spins about its axis of maximum moment of inertia, then the
attitude of the satellite is stable. By applying the Routh–Hurwitz criteria to this case,
we want to ascertain whether this attitude is also asymptotically stable.
The values of the coefficients of the fourth-degree polynomial are
a4 ¼ mð1 lÞIx Iy ¼ 10 ð1 0:01Þ 3002 ¼ 891000

a3 ¼ cIx ½Iy þ mR2 ð1 lÞ ¼ 150 300

300 þ 10 12 ð1 0:01Þ ¼ 13945500

a2 ¼ k½Iy þ mR2 ð1 lÞIx þ mð1 lÞ½ðIx Iz Þ Iy Iz mR2 ð1 lÞIx x20

¼ 10000 300 þ 10 12 ð1 0:01Þ 300 þ 10 ð1 0:01Þ
½ð300 500Þ2 10 12 ð1 0:01Þ 300 ð2 3:1416Þ2 ¼ 9:4417 108

a1 ¼ cf½Ix Iz mR2 ð1 lÞ Iy Iz gx20 ¼ 150 f½300 500 10 12
ð1 0:01Þ ð300 500Þg ð2 3:1416Þ2 ¼ 2:4860 108

a0 ¼ kf½Ix Iz mR2 ð1 lÞ Iy Iz gx20 þ ½ Iy Iz ð1 lÞ2 m2 R2 x40
¼ 10000 f½300 500 10 12 ð1 0:01Þ ð300 500Þg
ð2 3:1416Þ2 þ ½ð300 500Þ ð1 0:01Þ2 102 12
ð2 3:1416Þ4 ¼ 1:6542 1010
Since the values of all the coefficients computed above are greater than zero,
then the first condition is satisfied. As to the other conditions, there results
a3 a2 a4 a1 ¼ 13945500 9:4417 108

891000 2:4860 108 ¼ 1:2945 1016
a3 a2 a1 a4 a21 a23 a0 ¼ 13945500 9:4417 108
2
2:4860 108 891000 2:4860 108
ð13945500Þ2 1:6542 1010 ¼ 1:1934 1021
The positive values computed above show that all the conditions indicated above
are satisfied. Therefore, the attitude of the satellite is asymptotically stable.
Finally, the location illustrated in the preceding figures for a linear-tube damper
is not the only one possible. The following figure, due to the courtesy of NASA
[21], shows alternative locations for the same device.
In addition, the following figure, which is also due to the courtesy of NASA
[21], shows possible locations for a nutation damper of the fluid ring type.
Further types of passive nutation dampers (e.g. ball-in tube, cantilevered mass)
are described and illustrated in Ref. [49].
9.10 Gravity-Gradient Stabilisation of Satellites
The gravity gradient provides stabilising torques to a prolate satellite revolving

around the Earth in a circular orbit, when the principal axis of minimum moment of
inertia of the satellite is directed towards the centre of the Earth. In order for the
attitude of the satellite to be constant, the satellite must rotate about its axis per-
pendicular to the orbital plane at an angular velocity
l 12
xo ¼ E
r3
which expression follows from the vis-viva integral v2 = lE(2/r − 1/a). Only this
case is considered below, though steady motions of the satellite are also possible
about the same axis at angular velocity which differs from xo [50].
The attitude of the satellite in a circular orbit is described here by means of the
roll, pitch, and yaw angles Ux, Uy, and Uz, as has been done in Sect. 9.5. Here,
again, we consider the orbital system of reference xoyozo, having its origin in the
centre of mass O of the satellite, and such that the xo-axis is the velocity vector of
centre of mass of the satellite, the yo-axis is perpendicular to the orbital plane, and
the zo-axis points constantly towards the centre of the Earth.
The angles Ux, Uy, and Uz indicate the attitude of the satellite at any time t > t0,
that is, the orientation of the body-fixed system xyz with respect to the system
xoyozo. The values of these angles are zero at the initial time t0, when the satellite is
in a point of equilibrium.
At another time t > t0, the attitude of the satellite differs from its attitude of
equilibrium. Let Ux, Uy, and Uz be the angular deviations (supposed to be small) of
the attitude of the satellite at time t from the attitude of equilibrium.
As has been shown in Sect. 9.1, the rates of change of the angular deviations Ux,
Uy, and Uz with time are related to the angular velocities xx, xy, and xz about the
body-fixed axes x, y, and z as follows
9.10 Gravity-Gradient Stabilisation of Satellites 1179
xx ¼ Uz sin Uy þ Ux xo cos Uy sin Uz

xy ¼ Uz cos Uy sin Ux þ Uz cos Ux xo ðsin Ux sin Uy sin Uz þ cos Ux cos Uz Þ
xz ¼ Uz cos Uy cos Ux Uz sin Ux xo ðcos Ux sin Uy sin Uz sin Ux cos Uz Þ
Since the values of Ux, Uy, and Uz are small, these expressions reduce to
xx ¼ Ux xo Uz
xy ¼ Uy xo
xz ¼ Uz þ xo Ux
By differentiating these expressions with respect to time and taking account of

the constant value of xo, there results
xx ¼ Ux xo Uz

xy ¼ Uy
xz ¼ Uz þ xo Ux
Because of the angular deviations Ux, Uy, and Uz, the components of the vector
MO (gravity-gradient torque) are those of Sect. 9.5, that is,
3
MOx ¼ x2o Iz Iy sinð2Ux Þ cos2 Uy
2
3
MOy ¼ x2o ðIz Ix Þ sin 2Uy cos Ux
2
3
MOz ¼ x2o Ix Iy sin 2Uy sin Ux
2
Again, since the values of the angular deviations Ux, Uy, and Uz are small, the
preceding equations reduce to

MOx ¼ 3x2o Iz Iy Ux
MOy ¼ 3x2o ðIz Ix ÞUy
MOz ¼ 0
The attitude of the satellite is governed by the Euler equations of Sect. 7.6:


By setting MOx = 3x2o(Iz − Iy)Ux, MOy = 3x2o(Iz − Ix)Uy, and MOz = 0 into the
Euler equations written above, there results

3x2o Iz Iy Ux ¼ Ix xx Iy Iz xy xz
3x2o ðIz Ix ÞUy ¼ Iy xy ðIz Ix Þxz xx

0 ¼ Iz xz Ix Iy xx xy
By setting xx = U`x − xoUz, xy = U`y − xo, xz = U`z + xoUx, x`x = U``x − xoU`z,
x`y = U``y, and x`z = U``z + xoU`x into the preceding equations, there results

Ix Ux þ 4x2o Iy Iz Ux þ xo Iy Iz Ix Uz ¼ 0
Iy Uy þ 3x2o ðIx Iz ÞUy ¼ 0

Iz Uz þ x2o Iy Ix Uz xo Iy Iz Ix Ux ¼ 0
As a result of the second of these equations, the rotational motion of the satellite
about the yo-axis is only due to variations of Uy. Therefore, this motion is inde-
pendent of the rotational motions of the satellite about the other two axes xo and zo.
By contrast, the rotational motions of the satellite about the axes xo and zo are
coupled one with the other, as a result of the other two equations.
The second equation may be rewritten as follows

Ix Iz
Uy þ 3x2o Uy ¼ 0
Iy
In case of Ix being greater than Iz, the solution of this equation is an undamped
oscillation, as follows

Uy0
Uy ¼ Uy0 cosðxtÞ þ sinðxtÞ
x
where the angular frequency of this oscillation is

1
Ix Iz 2
x ¼ xo 3
Iy
Again, in case of Ix being greater than Iz, the torque due to the gradient of gravity
can be used to stabilise the satellite in its rotational motion about the yo-axis.
Let us consider now the stability of the satellite in its rotational motions about
the axes xo and zo. The solutions of the other two (the first and the third) differential
equations are of the following type
Ux ¼ AUx expðktÞ
Uz ¼ AUz expðktÞ
where k is a complex number. The real part of k must be negative, in order for these
solutions to be stable. By setting Ux = AUx exp(kt) and Uz = AUz exp(kt) into the
other two differential equations, there results

½Ix k2 þ 4x2o Iy Iz Ux þ xo Iy Iz Ix kUz ¼ 0

xo Iy Iz Ix kUx þ ½Iz k2 þ x2o Iy Ix Uz ¼ 0
The values of k making these two equations dependent are those which deter-
mine the stability. In order to find these values, we impose the condition det(A) = 0,
where A is the following 2 2 matrix
" #
Ix k2 þ 4x2o Iy Iz x o Iy Iz Ix k
A¼
xo Iy Iz Ix k Iz k2 þ x2o Iy Ix
The condition det(A) = 0 gives rise to the following algebraic equation
Ix Iz k4 þ x2o ðIx Iy þ 2Iy Iz 3Iz2 þ Iy2 þ 2Ix Iz Þk2 þ 4x4o ðIy Iz ÞðIy Ix Þ ¼ 0
which is a quadratic equation in k2. In order for this equation to have two real roots,
its discriminant must be positive, that is, the following condition must be satisfied
x4o ðIx Iy þ 2Iy Iz 3Iz2 þ Iy2 þ 2Ix Iz Þ2 16Ix Iz x4o ðIy Iz ÞðIy Ix Þ [ 0
In addition, in order for these real roots to be both of them negative, all the
coefficients of the quadratic equation in k2 written above must be positive, in
accordance with the Cartesian rule of signs. In other words, not only must the
discriminant of the preceding equation be positive, but also the following three
conditions must be satisfied
Ix Iz [ 0
x2o ðIx Iy þ 2Iy Iz 3Iz2 þ Iy2 þ 2Ix Iz Þ [ 0
4x4o ðIy Iz ÞðIy Ix Þ [ 0
When all of the four conditions indicated above are satisfied, then the equation
Ix Iz k4 þ x2o ðIx Iy þ 2Iy Iz 3Iz2 þ Iy2 þ 2Ix Iz Þk2 þ 4x4o ðIy Iz ÞðIy Ix Þ ¼ 0
has two pairs of complex conjugate values of k. These two pairs correspond to two
undamped oscillations about the axes xo and zo.
These four conditions must be satisfied in addition to the condition Ix > Iz, which
comes from the second of the three differential equations written above. The
condition Ix > Iz must be satisfied for a stable oscillation about the yo-axis.
For a given radius r of a circular orbit, the value of xo = (lE/r3)½ is determined.
Therefore, a satellite can be stabilised by means of the gravity-gradient torque,
when the moments of inertia Ix, Iy, and Iz about its principal axes x, y, and z satisfy
the conditions indicated above. For this purpose, as has been shown by Gerlach
[11], the axis (z) of minimum moment of inertia must be directed towards the centre
of the Earth, and the axis (y) of maximum moment of inertia must be perpendicular
to the orbital plane. The axis (x) of intermediate moment of inertia lies in the orbital
plane and is directed along the orbital velocity vector. Such is the case with the
Moon in its orbit around the Earth [11].
Another region of stability, smaller than the one described above, exists where Iz
is greater than Iy. In this region, the attitude of the satellite is statically unstable
about the x0-axis, as shown by the differential equation

Ix Ux þ 4x2o Iy Iz Ux þ xo Iy Iz Ix Uz ¼ 0
which requires Iy > Iz for stability. However, the stabilising effect of the rotation
about the yo-axis dominates in this region. This holds only with a satellite which
does not dissipate kinetic energy. A satellite which does can be stabilised by the
gravity-gradient torque only in the first of the two regions described above.
The stability described above means that the amplitudes of the oscillations of the
satellite with respect to a point of equilibrium do not increase with time. In order for
the attitude of the satellite to be not only stable but also asymptotically stable, these
oscillations must be damped by using appropriate means. Such means must produce
torques in the correct relationship of phase with the torques due to the gravity
gradient. As has been shown in Sect. 9.9, this can be obtained by attaching an
auxiliary mass to the satellite and providing a mechanism which dissipates kinetic
energy between the satellite and the auxiliary mass.
Gerlach [11] cites several devices meant to this purpose. One of them is a control
moment gyroscope (see Sect. 9.14). This device uses energy to spin the gyroscope,
but not to provide a torque on the satellite. Its primary function is to damp the
oscillations of the satellite, and its secondary function is to stiffen the roll and yaw
axes of the satellite.
The secondary mass attached to a satellite to damp its oscillations may also
consist of one or more rigid bodies. In this case, the dissipation of kinetic energy
occurs in a hinge connecting the satellite with the auxiliary bodies.
In another type of mechanism, a secondary mass is attached to the satellite
through a spring meant to dissipate kinetic energy and fastened to the end of a long
boom.
Still another mechanism uses two auxiliary bodies in the form of long rods
attached to the satellite by means of hinges. Both of these damping rods lie in the
horizontal plane in conditions of equilibrium. One of them is aligned with the
velocity vector of the satellite and damps the oscillations which occur about the axis
perpendicular to the orbital plane. The second damping rod is placed perpendicu-
larly to the first and is meant to damp the other two angular oscillations. A vertical
rod in the satellite increases the moments of inertia about the two principal axes
contained in the horizontal plane. When the orbit of a satellite is not exactly
circular, the zo-axis of the orbital system of reference xoyozo attached to the satellite
is not always directed towards the centre of the Earth. In this case, Koelle [23]
suggests to compute the angular velocity x of the radius vector going from the
centre of the Earth to the satellite as follows

5

x ¼ xo 1 þ 2e cos xo t tp þ e2 cos xo t tp þ
2
where xo = (lE/a3)½, lE is the gravitational parameter of the Earth, tp is the time of

passage at perigee, and a and e are, respectively, the major semi-axis and the
eccentricity of the orbit.
The eccentricity of the orbit causes periodic variations of the torque due to
gravity gradient, and hence oscillations of the satellite about the axis (yo) perpen-
dicular to the orbital plane. Therefore, the allowable eccentricity of the orbit is to be
considered carefully, in case of a satellite stabilised by means of the
gravity-gradient torque.
9.11 Stabilisation of Satellites by Means

of Electromagnetic Induction
The magnetic field due to the Earth can be used to generate torques which control
the attitude of a satellite. These torques are generated by placing either a permanent
magnet or an electromagnet along a desired axis, for the purpose of aligning one of
the principal axes of inertia of a satellite with the local direction of the magnetic
field of the Earth. This alignment is shown in the following figure, which is due to
the courtesy of the Michigan Exploration Laboratory of the University of Michigan
[52].
A magnetically controlled satellite behaves like the needle of a compass which

aligns with the magnetic dipole of the Earth.
The magnetic field of the Earth can be used not only to control the attitude of a
satellite but also to damp its angular oscillations about a desired position. This
damping is obtained by using rods of magnetic hysteresis material fixed to the
chassis of a satellite. By so doing, eddy currents are generated which damp the
oscillations about the centre of mass. The damping obtained by means of magnetic
hysteresis is based on the following principle. When a rod made of a ferromagnetic
material rotates in a magnetic field, a quantity of energy DEh is dissipated because
of the motion of the magnetic domains within the material. The energy loss over
any period of rotation of the rod is constant and is expressed, in joules/cycle, by
I
DEh ¼ V HdBi
where V is the volume of the rod in cubic metres, the integral in parentheses is the
area of the hysteresis loop, H is the intensity of the field in amperes/metre, and Bi is
the induced magnetic flux in the magnetic material. Since the energy loss per cycle
is independent of the spin rate, the spin rate decreases linearly with time and
reduces to zero with respect to the ambient field in a finite interval of time [45]. This
type of passive magnetic stabilisation has been used by the University of Michigan
for the M-Cubed satellite [52].
Care must be taken at the end of the powered flight of a space vehicle, in the
transition from the spin stabilisation to the magnetic stabilisation of the vehicle.
Electric currents flowing in coils placed within a satellite also produce torques,
because the magnetic fields associated with such currents interact with the magnetic
9.11 Stabilisation of Satellites by Means of Electromagnetic Induction 1185
field of the Earth. Torques generated in this way can be used to change the attitude
of a satellite with respect to the Earth. For example, the orientation of the TIROS 3
satellite was effectively controlled by applying currents in a loop to alter the
magnetic moment [43].
By the way, TIROS 3 (Television and Infrared Observation Satellite,
NSSDC/COSPAR ID 1961-017A) was a spin-stabilised meteorological satellite
designed to test experimental television techniques and infrared equipment. It is
shown in the following figure, due to the courtesy of NASA [38].
The spin axis of the TIROS 3 satellite was oriented with an accuracy of 1 or 2°
by means of a magnetic control device consisting of 250 cores of wire wound
around the outer surface of the satellite. The interaction between the magnetic field
induced in the satellite and the magnetic field of the Earth provided the necessary
torque for attitude control.
A device called magnetorquer (or magnetic torquer) has been proposed by
Kamm [22]. This device consists of three magnetometers and three current-carrying
coils placed along the principal axes of inertia of a satellite. Each magnetometer
measures the instantaneous flux density along its own axis. In order to avoid
interferences between the measurements and the magnetic torques generated, the
magnetometers and the magnetic torquers are used in different times. This type of
magnetic control on three axes has also been studied by White et al. [19].
A bar made of a permanent magnet mounted on gimbals can be used instead of
coils for the purpose of generating variable torques. This magnet can be used not
only to establish the direction of the external field but also to apply the magnetic
torque of the desired magnitude to the proper axis. This device has been described
by Adams and Brissenden [17]. A simulator of this device has also been described
by Adams et al. [18]. A scheme of this apparatus is shown in the following figure,
due to the courtesy of NASA [18].
A magnetic torquer mounted on gimbals has also been described by Griffin and
French [26].
Magnetic systems for attitude control provide a simple way to stabilise satellites
orbiting around a planet having a magnetic field, as is the case with the Earth. The
hardware used by them is simple, lightweight, and not subject to degradation or
variation of mass with time. They can work indefinitely, by using the electric power
coming from the solar panels to overcome the resistive loads due to the coils. They
are also reliable, due to the absence of moving parts.
By contrast, they have some disadvantages and limitations. A magnetic moment,
meant to control the attitude of a satellite, can only be applied perpendicularly to the
direction of the local magnetic field. Because of their dependence on the strength of
a planetary magnetic field, they are suitable for satellites orbiting at low altitudes,
but not for those orbiting at geosynchronous altitude. In addition, the magnetic field
of the Earth is highly variable. Also, the magnetic materials and the other pieces of
hardware of which a satellite is made can interact, and therefore both of them may
behave in undesirable ways.
9.12 Stabilisation of Satellites by Means of Reaction Jets
Reaction jets are used in discontinuous (pulsing) mode, specially for the purpose of
correcting errors in attitude and angular velocity at the moment in which a satellite,
after the powered flight, is injected into the desired orbit. For example, in case of a
satellite which has been spin-stabilised during the powered flight, reaction jets may
9.12 Stabilisation of Satellites by Means of Reaction Jets 1187
be used to de-spin the satellite before injection. They are also used to off-load
flywheels when external torques act predominantly in one direction. The principal
components of an attitude control system based on reaction jets of cold gas are
They are a storage tank for the gas (usually nitrogen or helium), which is
pressurised before launch, a pressure regulator, and a number of valves and nozzles
interconnected by means of pipes.
The gases named above are used because of their scarce or absent chemical
reactivity and low molecular mass. Helium has, over nitrogen, the advantage of a
higher (about 2.5 times) specific impulse, and therefore a significant reduction in
propellant mass, but at the expense of an increased (by a factor of about 7) storage
volume or pressure, and a higher cost [3].
In the scheme illustrated in the following figure, the attitude of a satellite is
controlled by a number of cold gas jets, each of which is parallel to one of the
principal axes x, y, and z of inertia of the satellite, and is placed at a common
distance ‘ from the centre of mass of the satellite.
Let ve be the exit velocity of the exhaust gas with respect to the satellite. Let m`x,
m`y, and m`z be the gas flow rates along the principal axes x, y, and z.
As has been shown in Sect. 4.2, in conditions of zero pressure at the exit from
the nozzles, the thrust forces along x, y, and z are respectively
fx ¼ mx ve
fy ¼ my ve
fz ¼ mz ve
In order for the motion of the centre of mass of the satellite not to be perturbed
by the forces due to the reaction jets, it is advisable to place always two parallel jets,
acting along opposite directions, on opposite sides of the satellite. By so doing, only
resulting moments and no resulting forces are generated by the cold gas jets. Since a
jet pipe can produce a force in only one direction, then 2 6 = 12 jet pipes are
necessary to control the attitude of a satellite shaped as a cube or a parallelepiped,
and then requiring two jet pipes on each face.
Because of the angular velocity vector x acting on centre of mass O of the
satellite, there are interactions between the linear velocities (ve) of the jets and the
angular velocities (xx, xy, and xz).
For example, with reference to the preceding figure, let us consider the gas jet
expelled from the nozzle indicated above by a grey arrow and pointing in the
positive direction of the z-axis. This jet has the following components of linear
velocity along the three body-fixed axes:
0 along the x-axis

‘xz along the y-axis
ve ‘xy along the z-axis
Therefore, this jet gives rise to the following contributions along the axes to the
reaction torque acting on the satellite:
0 along the x-axis

mz ðve ‘xy Þ‘ along the y-axis
mz ‘2 xz along the z-axis
By summing the contributions given by all the jets, the Euler equations shown in
Sect. 7.6


become, according to Ives [30] and Gerlach [11], those given below
9.12 Stabilisation of Satellites by Means of Reaction Jets 1189

MOx 2my ve ‘ 2‘2 ðmx þ my Þxx ¼ Ix xx Iy Iz xy xz
MOy 2mz ve ‘ 2‘2 ðmy þ mz Þxy ¼ Iy xy ðIz Ix Þxz xx

MOz 2mx ve ‘ 2‘2 ðmz þ mx Þxz ¼ Iz xz Ix Iy xx xy
The plus or the minus sign takes effect in the second terms on the left-hand side
of the preceding equations, depending on whether a positive or a negative torque is
desired. In case of a satellite whose jet pipes are arranged in a different manner, the
corresponding Euler equations are also different.
Since the value of the exhaust velocity ve is much greater than those of the linear
velocities ‘xx, ‘xy, and ‘xz, then the third terms on the left-hand side of the
preceding equations are negligible. Therefore, a proper placement of reaction jets
makes it possible to exert a torque about only one axis. In other words, the problem
of controlling the attitude of a satellite on three axes may be split up into three
separate problems.
The reaction jets used for controlling the attitude of a satellite are part of an on–
off control loop. The flow of gas through the jet duct is shut off by a control valve,
which can be either entirely open or entirely closed. The state of this valve depends
on the sign of a combination of a deviation in attitude and a rate of change of this
attitude. The following figure, redrawn from Gerlach [11], shows a typical attitude
control system using reaction jets.
In this figure, the deviation in attitude and its rate of change in time are denoted
by, respectively, U and U`. The transform variable is denoted by s. Such control
systems have also been studied by Abzug [1], Pistiner [40], Taylor and Smith [24],
and Brown [46].
According to Whitford [44], the magnitude of the thrust force due to reaction jets
ranges from 0.0005 to 0.5 N. More recently, Zandbergen [3] has given values
ranging from 0.001 to 266 N. The required magnitude of thrust may be determined
by either the desired angular acceleration of the satellite or the maximum perturbing
torque. In case of a satellite equipped with a system for orbital propulsion, Tripp
and Boardman [6] suggest the possibility of using thrust vector misalignment to
produce the maximum disturbance torque.
A disadvantage of the systems of stabilisation based on reaction jets is their high

mass consumption. Therefore, they are scarcely suitable for long missions.
9.13 Stabilisation of Satellites by Means of Reaction

Flywheels
A flywheel is a device for storing and regulating kinetic energy by means of a

continuously rotating mass, which is substantially a wheel mounted on an axle.
Such devices can provide energy storage and attitude control for satellites. When
two flywheels carried aboard a satellite spin in opposite directions at equal speeds,
then the satellite maintains its attitude. This is because, in the absence of torques,
the resultant moment is equal to zero. When an amount of kinetic energy of rotation
is transferred between the two flywheels to speed one of them and slow the other,
then the satellite rotates about its centre of mass, as a result of the imbalance of
moments. The amount of stored energy depends on the moment of inertia and the
rotational speed of a flywheel.
Flywheels have been used in the stabilisation systems of many satellites. For
example, the Nimbus 1 (NSSDC/COSPAR ID: 1964-052A), launched on 28
August 1964, was the first meteorological satellite stabilised on three axes. For this
purpose, it used gas jets and flywheels controlled by horizon sensors.
The satellite rotated once per orbit, by placing torques on the appropriate fly-
wheel, so that the instruments carried aboard the satellite pointed constantly
towards the Earth [47]. The following figure, due to the courtesy of NASA [37],
shows the Nimbus 1 satellite.
9.13 Stabilisation of Satellites by Means of Reaction Flywheels 1191
The method for controlling the attitude of a satellite by using reaction flywheels
is illustrated in the following figure.
The satellite shown in this scheme has three flywheels. The spin axes of the
flywheels are directed along the principal axes of inertia x, y, and z of the satellite.
The principal axes of inertia are assumed to maintain a fixed orientation in the
satellite, so that any fluctuations of the axes due to changes in distribution of matter
within the satellite are negligible. The flywheels are mounted in fixed bearings and
provide reaction torques when accelerated with respect to the body of the satellite.
The angular motions of roll, pitch, and yaw from the desired attitude are counter-
acted by the three flywheels indicated above, one along each principal axis. For
example, the flywheel along the roll axis of the satellite, when accelerated in the
same direction as the rolling motion of the satellite, generates an equal and opposite
reaction torque which is exerted on the satellite, so as to counteract the roll motion.
Let xx, xy, and xz be the three components of the velocity vector x along the
principal axes of inertia of the satellite. Let I be the moment of inertia common to
all flywheels along their respective axes. Let -x, -y, and -z be the angular
velocities of the flywheels with respect to the body of the satellite. The Euler
equations shown in Sect. 7.6


are to be modified in order to take account of the presence of the three flywheels.
The action of the flywheels does not change the total moment of momentum of
the combination comprising the body of the satellite and the flywheels. There is
only a transfer of moment of momentum between the body of the satellite and the
flywheels.
The Euler equations written above are modified to include the effects of the
flywheels. According to Ives [30] and Gerlach [11], there results

MOx ¼ Ix xx Iy Iz xy xz þ Ið-x þ -z xy -y xz Þ
MOy ¼ Iy xy ðIz Ix Þxz xx þ Ið-y þ -x xz -z xx Þ

MOz ¼ Iz xz Ix Iy xx xy þ Ið-z þ -y xx -x xy Þ
In the preceding equations, MOx, MOy, and MOz are the moments due to the
external torques acting to the satellite. Account is taken in Ix, Iy, and Iz of the
additional mass due to the flywheels. The ratio of the moments of inertia of the
flywheels to those of the body of the satellite is of the order of magnitude of 10−5
[11]
As shown by the preceding equations, the effects due to the external torques can
be counteracted by changing appropriately the angular acceleration of one or more
flywheels. In other words, the flywheels can isolate a satellite from perturbing
torques. In order for a satellite to be stable in a desired attitude, the acceleration of
its flywheels is made to depend on the angular deviations (U) and angular deviation
rates (U`) from this attitude. For this purpose, a simplified block diagram of a
control system on a single axis by means of a flywheel is shown in the following
figure, which is redrawn from Gerlach [11].
Flywheel control systems have also been studied by Abzug [1, 28], White and
Hansen [20], Cannon [42], and Roithmayr et al. [5]. As to the symbols used in the
preceding figure, s1 is the time constant associated with time-of-sight inputs, s2 is
the time constant associated with torque inputs, s is the transform variable, Km is the
motor gain, and sm is the time constant relating to motor and wheel.
The Euler equations written above, modified to take account of the effects due to
the three flywheels, show that a change in angular velocity of a flywheel rotating
about one of the three axes may have influence on the motion of the flywheels
rotating about the other axes. In order to avoid this disadvantage, Ormsby and
Smith [39] and Hering and Hufnagel [14] have proposed to use a single reaction
sphere instead of three reaction wheels. In particular, Hering and Hufnagel have
9.13 Stabilisation of Satellites by Means of Reaction Flywheels 1193
shown that an electrically conducting sphere placed in a rotating magnetic field is

subject to a torque about its axis of rotation. This torque, which is due to eddy
current drag between the field and the sphere, can be used to control the attitude of a
satellite. Angular momentum accumulated in the satellite appears as rotation of the
sphere and is automatically dumped by interaction of the eddy current with the
magnetic field of the Earth. This sphere can be placed without mechanical contact
with the satellite by high-frequency alternating magnetic fields. This type of control
on three axes by using one sphere eliminates gyroscopic interaction [14]. In
mathematical terms, the cross-coupling terms, which would otherwise be present in
the Euler equations, are eliminated [11].
An external torque acting constantly on a satellite in the same direction produces
the effect of accelerating continuously at least one of its flywheels. For practical
reasons, the maximum allowable angular velocity of the flywheels cannot exceed a
certain value. When this value is reached, a flywheel becomes saturated. Therefore,
a satellite which is expected to experience torques acting in only one direction must
be provided with an additional mechanism for attitude control. For this purpose,
permanent magnets or current-carrying coils may be used, so as to generate torques
through interaction of their own magnetic fields with the magnetic field of the Earth,
as has been shown in Sect. 9.11. Adams and Brissenden [17] and Adams et al. [18]
have described this additional mechanism. Reaction jets, described in Sect. 9.12,
may be used for the same purpose, as has been proposed by Haeussermann [13]. By
using these reaction jets when required, the flywheels can be desaturated, because
their angular velocities are brought to values near to zero.
The following figure, redrawn from Gerlach [11], shows a block diagram
relating to an attitude control working in dual mode on a single axis.
A system of attitude control based on flywheels is appropriate to counteract

torques whose sign changes periodically, as is the case with several perturbing
torques which act on satellites orbiting around the Earth.
9.14 Stabilisation of Satellites by Means of Control

Moment Gyroscopes
The reaction flywheels described in Sect. 9.13 have their spin axes constantly
aligned with the principal axes of inertia x, y, and z of the body of the satellite which
carries them. They counteract the perturbing torques acting on the satellite by only
increasing or decreasing the magnitude of their angular speeds, which vary from
zero to several thousands of rotations per minute [8].
By contrast, a control moment gyroscope is a torque actuator based on a wheel
rotating at a constant speed and mounted on either a single gimbal or a set of
gimbals. By applying torques to the gimbal(s), the orientation of the wheel with
respect to the axes x, y, and z is changed, and therefore the direction of the
moment-of-momentum vector of the wheel is also changed.
Control moments gyroscopes are available as single- or double-gimballed
devices [15]. One of the four double-gimballed control moment gyroscopes which
are mounted on the International Space Station is shown in the following figure, due
to the courtesy of NASA [4].
As shown in this figure, a control moment gyroscope has the following essential
parts:
9.14 Stabilisation of Satellites by Means of Control Moment Gyroscopes 1195
• a wheel rotating about its spin axis at a large and usually constant speed, in order
to produce a moment-of-momentum vector of constant magnitude;
• a spin motor acting on the wheel;
• a single or double gimbal, on which the wheel is mounted; and
• a gimbal motor, or a set of gimbal motors, generating torque(s) for the purpose
of changing the moment-of-momentum vector of the wheel to the desired
direction.
The torque exerted by the motor on the gimbal generates a motion of precession
directed perpendicularly to the gimbal axis and to the spin axis of the wheel.
A small input torque to the control moment gyroscope generates a large output
torque to the satellite. This fact makes control moment gyroscopes particularly apt
to generate large control torques and also to store large moments of momentum
over long periods of time.
Control moments gyroscopes can be classified as follows:
• single-gimballed devices, wherein the wheel is gimballed along one axis and
constrained to rotate in a plane perpendicular to this axis;
• double-gimballed devices, wherein the wheel is gimballed along two perpen-
dicular axes, and the moment-of-momentum vector of constant magnitude is
contained in a sphere;
• variable-speed devices, wherein the wheel of a single-gimballed device rotates
at a variable speed.
Single-gimballed devices are cheap, lightweight, and reliable, but are subject to
singularities. A singularity is the impossibility of generating torques about a certain
direction, and therefore the loss of control of a spacecraft on three axes.
The dynamic behaviour of a control moment gyroscope is governed by the
following equations. With reference to the following figure, let n be the spin axis of
the wheel.
A control moment gyroscope is similar to a common gyroscope, but operates in

the opposite way. When the spinning wheel of a control moment gyroscope is made
to rotate about its gimbal axis f at a small precession rate d` by an input torque, then
a large rotation, and therefore an output torque NO, arises about the η-axis.
The left-hand side of the preceding figure shows a control moment gyroscope in
its initial position (at time t = 0). Let = be the matrix of inertia of the wheel.
Let hO = = x be the moment-of-momentum vector applied to the rotating
wheel at t = 0 along its spin axis n. After an interval Dt of time, the wheel is rotated
about its gimbal axis f by an angle d at a rate d` by an input torque. This rotation
gives rise to an output torque NO about the η-axis. The torque NO may be con-
sidered as resulting from two forces f, equal in magnitude and oppositely directed,
which act in the points A and B placed along the n-axis.
The rotation of the wheel about its gimbal axis f causes the
moment-of-momentum vector hO of the wheel to also rotate by an angle Dd. Let
h*O be the moment-of-momentum vector of the wheel after the time interval Dt. Let
DhO be the change in moment of momentum in the interval Dt. The vector DhO is
parallel and oppositely directed with respect to the angular impulse vector NODt.
By inspection of the preceding figure, the magnitude of the vector DhO results
jDhO j ¼ hO Dd
By equating the magnitude of the angular impulse to the magnitude of the

change in moment of momentum, there results
NO Dt ¼ hO Dd
By dividing all terms of the preceding equation by Dt and letting Dt tend to zero,
there results
dd
NO ¼ hO hO d
dt
where hO = Ix, I is the moment of inertia of the wheel about its spin axis n, x is the
spin rate of the wheel, and d` is the gimbal rate. The preceding equation, rewritten
in vector form, expresses the output torque as follows
N O ¼ hO d
where d` is the gimbal rate vector. The attitude of a satellite can be controlled on
three axes by means of a cluster of control moment gyroscopes, as will be shown
below. The rotational motion of a satellite equipped with control moment gyro-
scopes is governed by the following equation
Hs þ x H s ¼ N ext
where Hs is the moment of momentum with respect to the control axis of the
satellite in the system of reference xyz attached to the body of the satellite and
having its origin in the centre of mass G of this body, Next is the vector of the
external torques, which includes all types of perturbing torques, and x is the
angular velocity vector of the satellite with respect to an inertial system.
The total moment of momentum of the satellite is
H s ¼ =s x þ h
where =s is the matrix of inertia of the satellite, and h is the moment-of-momentum

vector associated with the cluster of control moment gyroscopes.
By substituting
H s ¼ =s x þ h
into
Hs þ x Hs ¼ N ext
there results
=s xþ hþ x ð=
=s x þ hÞ ¼ =s xþ sþ x ð=
=s xÞ þ x h ¼ N ext
Let U be the internal control torque generated by the cluster of control moment
gyroscopes. When we introduce U in the preceding equation, the torque resulting
from the sum of the two terms h` and x h must be equal to −U. Therefore
hþ x h ¼ U
=s xþ x ð=
=s xÞ ¼ N ext þ U
Let us assume the internal control torque U to be known and chosen in such a
way as to satisfy the equation
hþ x h ¼ U
Let n be the number of the single-gimballed control moment gyroscopes which

form the cluster. The moment-of-momentum vector h is a function of the gimbal
angles d1, d2, …, dn, each of which relates to one of the n control moment gyro-
scopes. As has been shown above, a cluster of three constant-speed,
single-gimballed control moment gyroscopes is subject to singularities at certain
gimbal angles. Let us consider a cluster comprising four constant-speed,
single-gimballed control moment gyroscopes, as shown in the following figure,
which is redrawn from Ref. [53].
This case has been studied by Bedrossian [31] and then by Lappas [53], and is
also illustrated in the following figure, due to the courtesy of NASA [31], where
uf1, uf2, uf3, and uf4 are unit vectors along the gimbal axes of the four control
moment gyroscopes, respectively, 1, 2, 3, and 4, which form the cluster.
As shown on the left-hand side of the preceding figure, this cluster of four
constant-speed, single-gimballed control moment gyroscopes is arranged to form a
pyramid, whose skew angle b is equal to 54°.73. According to Lappas [53], this
pyramid mounting arrangement is the best in terms of the uniformity of the

envelope of the moment-of-momentum vector.
By the way, a moment-of-momentum envelope, shown on the right-hand side of
the preceding figure, is a three-dimensional surface defined as the locus of all points
traced out by the maximum moment of momentum attainable by a cluster of control
moment gyroscopes along any given direction. In the case of the pyramid mounting
arrangement illustrated above having a skew angle b = 54°.73, this surface has a
nearly spherical shape, except for funnel-shaped cavities at certain directions, which
represent states of unattainable moment of momentum, because the normal to the
surface passing through each funnel is aligned with the gimbal axis [31]. Therefore,
apart from these directions, the moment-of-momentum vector hi (i = 1, 2, 3, and 4)
has about the same magnitude for all the three orthogonal axes x, y, and z of the
body-fixed system of reference.
By projecting the four moment-of-momentum vectors onto the axes x, y, and z,
Bedrossian [31] has determined their components for each of the four control
moment gyroscopes in the pyramid mounting arrangement indicated above, as
follows
2 3 2 3
cos b sin d1 cos d2
6 7 6 7
h1 ¼ hO 4 cos d1 5 h2 ¼ hO 4 cos b sin d2 5
sin b sin d1 sin b sin d2
2 3 2 3
cos b sin d3 cos d4
6 7 6 7
h3 ¼ hO 4 cos d3 5 h4 ¼ hO 4 cos b sin d4 5
sin b sin d3 sin b sin d4
where hO is the constant magnitude of the moment-of-momentum vector of each

spinning wheel. The total moment-of-momentum vector due to the cluster of four
control moment gyroscopes is
h ¼ h1 þ h 2 þ h 3 þ h 4
The output torque of the cluster results from differentiating the preceding
equation with respect to time in the body-fixed system of reference:
h ¼ h1 þ h2 þ h3 þ h4 ¼ JðdÞd
where J(d) is the 3 4 Jacobian matrix whose entries are the partial derivatives of
the total moment-of-momentum vector h due to the cluster of moment control
gyroscopes with respect to the gimbal angles dj (j = 1, 2, 3, and 4), as follows
2 3
cos b cos d1 sin d2 cos b cos d3 sin d4
@hi
Jð dÞ ¼ ¼ hO 4 sin d1 cos b cos d2 sin d3 cos b cos d4 5
@dj
sin b cos d1 sin b cos d2 sin b cos d3 sin b cos d4
and d` is the 4 1 vector containing the rates of change with time of the gimbal
angles d1, d2, d3, and d4 of the four control moment gyroscopes.
As mentioned above, the value of the internal control torque U is known and is
chosen in such a way as to satisfy the equation h` + x h = −U. Therefore, h` is
determined.
A possible manner (called pseudo-inverse steering law) to determine the d`
vector is to use the pseudo-inverse, J+, of the Jacobian matrix J, as follows
1
d ¼ J þ h ¼ JT JJT h
A method for computing the pseudo-inverse, A+, of a given nonzero

n m matrix A has been shown at length in Sect. 2.10. According to this method,
called singular value decomposition, the 3 4 Jacobian matrix J may be
decomposed by using its singular values d1, d2, and d3, as follows
J ¼ PDQT
where P is a 3 3 matrix with orthonormal columns (such that PTP = I), QT is a

4 4 matrix with orthonormal columns (such that QTQ = I), and D is the fol-
lowing 3 4 diagonal matrix whose nonzero entries are the singular values d1, d2,
and d3 of the Jacobian matrix J, as follows
2 3
d1 0 0 0
D¼4 0 d2 0 05
0 0 d3 0
The 4 3 pseudo-inverse matrix, J+, of the Jacobian matrix J results from
J þ ¼ QD þ PT
where D+ is the following 4 3 diagonal matrix

2 3
1=d1 0 0
6 0 1=d2 0 7
Dþ ¼6
4 0
7
0 1=d3 5
0 0 0
whose nonzero entries are the reciprocals of the singular values of the Jacobian
matrix J. When the singular values of J have been determined, the steering law
based on the pseudo-inverse matrix can be expressed as follows
" #
3
X 1
d ¼ J þ h ¼ qi pTi h
i¼1
di
where pi and qi are the column vectors which form the matrices, respectively, P and
Q of the singular value decomposition J = PDQT.
However, when the rank of the Jacobian matrix J is less than three, then the
pseudo-inverse matrix J+ does not exist, and therefore the control system has a
singularity. In other words, the steering law based on the pseudo-inverse matrix
fails when the rank of J is less than three.
A measure of the rank of J is given by the following quantity

1
m ¼ det JJT 2
where |det(JJT)| is the absolute value of the determinant of the 3 3 matrix

resulting from the product JJT. When the value of m approaches zero, then the
cluster of control moment gyroscopes described above cannot control anymore the
attitude of a satellite along one of the three body-fixed axes.
In order to avoid this, other steering laws than that based on the pseudo-inverse
matrix J+ have been proposed by several authors. Some of them are briefly
described below.
One of these methods, known as singularity-robust steering law, is based on a
mathematical artifice, which alters the 3 3 matrix resulting from the product JJT
by adding a perturbing matrix (that is, kI) to it, according to the following equation
h 1 i
d ¼ J# h ¼ JT JJT þ kI h
where k is a positive scalar quantity, and I is the 3 3 identity matrix, as follows

2 3
1 0 0
I ¼ 40 1 05
0 0 1
Nakamura and Hanafusa [56], to whom this method is due, suggest to select the
value of k as follows
k¼0 for m m0
k ¼ k0 ð1 m=m0 Þ2 for m\m0
where k0 and m0 are two constant quantities, whose values are to be selected by the
designer. Other expressions of k are also possible, as will be shown below.
Remembering the singular value decomposition J = PDQT, it is possible to
write another version of the singularity-robust steering law, as follows
" #
X3
di
d ¼ q pT h
i¼1
di2 þ k i i
The steering law written above reduces to the steering law of the pseudo-inverse
method, when k = 0.
Another method, known as modified (or general) singularity-robust steering law,
also adds a perturbation to the 3 3 matrix resulting from the product JJT, but this
time the steering law is governed by the following equation
h 1 i
d ¼ J# h ¼ JT JJT þ kE h
where E is the following 3 3 symmetric and positive-definite matrix

2 3
1 3 2
E ¼ 4 3 1 1 5
2 1 1
where e1, e2, and e3 are periodic functions of time having the following form
ei ¼ e0 sinðxt þ ui Þ
(i = 1, 2, and 3). In the preceding equations, k = k0 exp[−l det(JJT)], e0 is the

amplitude, and x and ui are, respectively, the angular frequency and the phase
offsets of modulation of the three periodic functions ei. The values of the quantities
k0, l, e0, x, u1, u2, and u3 are to be chosen by the designer. This method, which is
due to Wie et al. [2], uses the dither signals ei in order to generate the quantity kE,
which perturbs the product JJT.
Still another modified singularity-robust steering law has been suggested by
Okubo and Tani [12]. They have defined a degree of singularity as follows
d1
j
d3
where d1 and d3 are, respectively, the maximum and the minimum of the singular
values of the Jacobian matrix J. The value of j increases when the values of the
gimbal angles d1, d2, d3, and d4 are close to a singular state, because the value of the
minimum singular value d3 becomes very small. The singularity-robust steering law
proposed by Okubo and Tani [12] can be expressed as follows
" #
X3
di 1 j j0
d ¼ q p T
hþ k SA þ 1 q1
i¼1
di2 þ k i i 2 jj j0 j
where k, kSA, and j0 are quantities whose values are to be chosen by the designer.
The right-hand side of the preceding equation results from the sum of two terms.
The first term is the same as the only term which appears in the steering law
proposed by Nakamura and Hanafusa [56], which has been described above.
The second term has been added to the first by Okubo and Tani [12] in order for
the control system to escape rapidly from a singular state, by giving a gimbal rate
command which provides a torque in the direction orthogonal to the singular
direction. When j is smaller than j0, then the second term is equal to zero, and the
steering law is the same as that of Ref. [56]. The singular vector q1 is used in the
second term in order to exert the maximum possible torque in the direction
orthogonal to the singularity surface, and to escape rapidly from a singular point.
In case of an inadequate steering law being chosen, the control system of a
satellite may go into a singular state from which it cannot escape. In this case, an
external torque is required in order for the control system to leave this state. This
torque may be provided by a different actuator, such as a flywheel, a magnetic
torquer, or a thruster pair. By applying this torque, a cluster of single-gimballed
control moment gyroscopes returns to a fully controllable state.
Otherwise, double-gimballed constant-speed or single-gimballed variable-speed
control moment gyroscopes may be used.
9.15 Three-Axis Controlled Satellites
The methods shown in Sects. 9.12 (mass expulsion), 9.13 (reaction flywheels), and
9.14 (control moment gyroscopes) can be used, either alone or combined, to control
actively a spacecraft on three axes. For example, each of the Voyager 1 and 2
spacecraft was equipped with 16 hydrazine thrusters to be used only for attitude
control. Since 100 kg of propellant are carried on board the spacecraft for this
purpose and the consumption is about 1.6 g/day, then this quantity of propellant
suffices for about 0.0016 365.25 100 = 58.44 years.
As another example, the Cassini–Huygens spacecraft, launched on 15 October
1997 and directed towards Saturn and Titan, was equipped with 16 0.5 N hydrazine
thrusters and also with 3 main reaction wheel assemblies for attitude control, as
shown in the following figure, due to the courtesy of NASA-JPL [34].
By contrast, a spin-stabilised spacecraft rotates around its own vertical axis and
spins like a top. This keeps the attitude of the spacecraft under control.
Three-axis stabilisation has advantages and disadvantages with respect to spin
stabilisation. A spinning cylindrical satellite is stabilised by gyroscopic effect,
which creates an inertial stiffness and keeps its spin axis in the desired direction.
This is advantageous for fields and particles instruments, but requires complex
instruments to de-spin antennas or optical instruments which must be pointed at
fixed targets for scientific observation of celestial bodies or communications
between the satellite and the Earth. In addition, a spin-stabilised satellite cannot use
large solar panels to obtain power from the Sun. Therefore, it requires large
amounts of battery power.
A three-axis stabilised satellite can point its antennas and optical instruments at
the desired direction without the necessity of de-spinning them. It can also point its
solar panels at the desired direction with respect to the Sun. However, it may have
to execute special manoeuvres, which are required for the best use of its fields and
particles instruments.
When thrusters are used for routine stabilisation, then a three-axis stabilised
satellite rocks slowly back and forth within a dead band of allowed attitude error,
9.15 Three-Axis Controlled Satellites 1205
that is, within a range of allowed positions. This motion is not always predictable,
which is a disadvantage when a satellite carries observation instruments. A satellite
stabilised by reaction flywheels is much steadier than one stabilised by thrusters, but
has a greater mass. In addition, as has been shown in Sect. 9.13, reaction flywheels
have a limited mechanical lifetime and require frequent desaturation manoeuvres.
9.16 Attitude Re-orientation of a Satellite by Means

of Impulse Coning
The manoeuvre described in the present section is called impulse coning. It is a

method of attitude re-orientation, in which impulsive torques are used to first ini-
tiate and then terminate a free motion of precession of a spinning satellite. Between
a torque impulse and the following one, the satellite goes undisturbed towards the
desired attitude. These impulsive torques are applied to the satellite by means of
onboard thrusters.
With reference to the following figure, let G and hG0 be, respectively, the centre
of mass and the moment-of-momentum vector, at the initial time t = 0, of a
cylindrical satellite spinning about its longitudinal axis of symmetry. It is desired to
execute an impulse coning manoeuvre, which changes only the direction but not the
magnitude of the vector hG0, so that this vector rotates in space through an angle h.
For this purpose, it is necessary to apply an external moment MG such that the
resulting variation DhG of the moment-of-momentum vector, at t = Dt, is
ZDt
DhG ¼ M G dt
0
The external moment MG is applied to the satellite perpendicularly to its spin

axis by means of thrusters. These thrusters exert two equal and oppositely directed
forces f. As a result of this external moment, the spin axis z of the satellite has a
motion of precession about an axis which forms a constant angle with the
moment-of-momentum vector hG0. Let h/2 be this constant angle.
As has been shown in Sect. 7.8, the spin axis z of this cylindrical (such that
Ix = Iy) satellite rotates about the vector hG0 at a precession rate w` which is
proportional to the spin rate /` of the satellite. Therefore, the following equality
holds
Iz /
w ¼
h
ðIz Ix Þ cos
2
where w, h, and / are the classical Eulerian angles defined in Sect. 7.8, and x, y, and
z are the principal axes of inertia of the satellite.
Since the satellite shown in the preceding figure is a prolate (Ix > Iz) rigid body,
then the preceding equality can be rewritten as follows
Iz /
w ¼
Ix Iz h
cos
2
As a result of the first change (DhG1) of moment of momentum, the spin axis of
the satellite begins its motion of precession about a fixed axis. When the tip of the
moment-of-momentum vector hG has travelled through a semi-circumference, then
a further external moment MG is applied to the satellite. This causes another change
9.16 Attitude Re-orientation of a Satellite by Means of Impulse Coning 1207
(DhG2) of moment of momentum perpendicularly to the spin axis and in the same
direction, with respect to the satellite, as the previous change (DhG1), so that the
following equality holds
jjDhG2 jj ¼ jjDhG1 jj
This, as indicated above, terminates the free motion of precession of the spin
axis and makes the satellite stable in its new attitude.
The time t1 necessary to execute a single manoeuvre aimed at re-orienting the
spin axis of a satellite through an angle h is found by dividing the path (a
semi-circumference or p radians) travelled by the tip of the vector hG by the
precession rate w`, as follows

p pðIx Iz Þ h
t1 ¼ ¼ cos
w /Iz 2
The expenditure of propellant depends on the magnitude of each single increment

||DhG|| of moment of momentum. Such is also the case with the increments Dv of
velocity necessary to execute the orbital manoeuvres described in Chaps. 4 and 5.
The preceding figure shows that the total increment of moment of momentum
necessary for the purpose of re-orienting the spin axis of a satellite through an angle
h in a single manoeuvre is
jjDhG jj ¼ jjDhG1 jj þ jjDhG2 jj ¼ 2jjhG0 jj tanðh=2Þ
where ||hG0|| is the magnitude of the moment-of-momentum vector hG0 of the

satellite at the initial time t = 0.
The same result can also be obtained by means of a sequence of n impulse
coning manoeuvres, each of which re-orients the spin axis of a satellite by a small
angle a such that na = h. Since a is a small angle, then tan(a/2) a/2, and
therefore in this case
jjDhG jj 2njjhG0 jja=2 ¼ jjhG0 jjna ¼ jjhG0 jjh
For example, in case of the desired angle of rotation being h = 30° = p/6
radians, a single impulse coning manoeuvre requires a total increment of moment of
momentum ||DhG|| = 2||hG0||tan(p/6), whereas a sequence of ten impulse coning
manoeuvres, each of which changes the spin axis by 3° = p/60 radians, requires a
total increment of moment of momentum ||DhG|| = ||hG0||p/6.
Since 2 tan(p/6) 1.155 and p/6 0.5236, then the second method is better
than the first, because it requires a smaller total increment of moment of momen-
tum. As is easy to verify, this advantage increases for higher values (40°, 50°, …) of
the desired angle of rotation.
On the other hand, the time tn taken by a sequence of several small manoeuvres
is higher than the time t1 taken by a single manoeuvre. In addition, a single
manoeuvre is safer than a sequence of small manoeuvres.
Remembering the preceding equation

pðIx Iz Þ h
t1 ¼ cos
/Iz 2
and that a = h/n, there results

pðIx Iz Þ 1h
tn ¼ n cos
/Iz 2n
Therefore, the ratio of tn to t1 is

h
cos
tn 2n
¼n
t1 h
cos
2
For example, if h = 30° and n = 10, then

cos 1
:5
tn ¼ 10 t1 ¼ 10:349 t1
cos 15
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Chapter 10
Dynamics of Spinning Rockets
10.1 The Motion of a Spinning Rocket
The present chapter deals with the attitude and the motion of the centre of mass of a
spinning rocket subject to thrust. The mass of the rocket is considered at first
constant in time, and then variable as a result of the thrust applied. A rocket is
considered as a rigid body which rotates at an angular velocity x with respect to an
inertial system of reference XYZ and has a body-fixed system of reference xyz,
whose origin O coincides with the centre of mass of the body. The axes x, y, and
z of the body-fixed system do not necessarily coincide with the principal axes of
inertia 1, 2, and 3 of the body itself. Therefore, it is necessary to consider not only
the moments of inertia Ix, Iy, and Iz of the body, but also its products of inertia Ixy,
Ixz, and Iyz.
As has been shown in Sect. 7.2, the components along x, y, and z of the
moment-of-momentum vector hO ¼ = x of a rigid body with respect to its centre
of mass O can be expressed as follows
hOx ¼ Ix xx Ixy xy Ixz xz

hOy ¼ Ixy xx þ Iy xy Iyz xz
hOz ¼ Ixz xx Iyz xy þ Iz xz
The fundamental equation M O ¼ hO þ x hO governing the motion of a rigid

body rotating about its centre of mass O, projected onto the axes x, y, and z, is

MOx ¼ Ix xx Ixy xy Ixz xz þ Ixz xx Iyz xy þ Iz xz xy Ixy xx þ Iy xy Iyz xz xz

MOy ¼ Ixy xx þ Iy xy Iyz xz Ixz xx Iyz xy þ Iz xz xx þ Ix xx Ixy xy Ixz xz xz

MOz ¼ Ixz xx Iyz xy þ Iz xz þ Ixy xx þ Iy xy Iyz xz xx Ix xx Ixy xy Ixz xz xy

https://doi.org/10.1007/978-3-319-62220-0_10
1214 10 Dynamics of Spinning Rockets
Let v and f be, respectively, the velocity vector of the centre of mass O and the
resultant force acting upon the body. The components of these vectors along the
axes x, y, and z of the body-fixed system of reference are respectively
v ¼ vx ux þ vy uy þ vz uz
f ¼ fx ux þ fy uy þ fz uz
Since the axes x, y, and z of the body-fixed system rotate together with the body
at an angular velocity x, then the equation of motion of the body is
f ¼ mv0 ¼ mv þ x mv
The preceding equation, projected onto the axes x, y, and z, yields the following
three equations
fx ¼ mðvx þ xy vz xz vy Þ
fy ¼ mðvy þ xz vx xx vz Þ
fz ¼ mðvz þ xx vy xy vx Þ
When the resulting force f does not pass through the centre of mass O of the
body, then the three equations expressing fx, fy, and fz become coupled with the
three equations expressing MOx, MOy, and MOz. In addition, these equations
describe the motion of the body only in terms of linear and angular velocities
relating to the body-fixed axes x, y, and z. The transformation of these equations in
terms of displacement and angles relating to an inertial system of reference XYZ is a
difficult task, which requires some simplifying assumptions to be carried out.
10.2 Misalignment of the Thrust Vector in Body-Fixed

Co-ordinates
With reference to the following figure, let us consider the problem concerning a
missile subject to a thrust force f which is misaligned with respect to the axis of
symmetry of the missile.
10.2 Misalignment of the Thrust Vector in Body-Fixed Co-ordinates 1215
Due to the symmetric shape of the missile, the axes x, y, and z of the system of
reference fixed to the body coincide with its principal axes of inertia 1, 2, and 3, and
also I1 ¼ I2 6¼ I3 . Since I1 = I2, it is possible to rotate arbitrarily the axes 1 and 2 so
as to make one of these axes (e.g. the axis 1) perpendicular to the plane which
contains the direction of the thrust force and the axis of symmetry 3.
Euler’s equations of Sect. 7.6, are rewritten below for convenience.

These equations, in the particular case considered here ðI1 ¼ I2 ; MO2 ¼ MO3 ¼ 0Þ,
reduce to

0 ¼ I2 x2 ðI3 I1 Þx3 x1
0 ¼ I3 x3
As a result of the third equation, x3 ¼ 0, and therefore x3 = constant. By setting
I1 I3
k ¼ x3
I1
the first two equations may be rewritten as follows
MO1 =I1 ¼ x1 þ kx2

0 ¼ x2 kx1
By multiplying all terms of the second equation by i ¼ ð1Þ1=2 and adding these
terms to the corresponding terms of the first equation, there results
MO1 =I1 ¼ x1 þ ix2 þ kðx2 ix1 Þ
Remembering that 1/i = −i, and setting for convenience x12 ¼ x1 þ ix2 , the
preceding equation may be written as follows
MO1 =I1 ¼ x1 þ ix2 ikðx1 þ ix2 Þ ¼ x12 ikx12
Thomson [1] solves the preceding differential equation by using the method of
Laplace transformation and finds the following equation expressed in terms of the
subsidiary variable s
x12 ð0Þ M O1 ðsÞ

x12 ðsÞ ¼ þ
s ik I1 ðs ikÞ
where x12 ðsÞ denotes the Laplace transform of x12(t).

The solution of this equation in the domain of time is [1]:
Zt
1
x12 ðtÞ ¼ x12 ð0Þe þikt
M O1 ðsÞeikðtsÞ ds
I1
0
The components x1 and x2 of the angular velocity vector x result from the
preceding equation by noting the real part and the imaginary part of x12.
10.3 Misalignment of the Thrust Vector in Inertial

Co-ordinates
In Sect. 10.2, the problem concerning the misalignment of the thrust vector has
been solved in the body-fixed system of reference xyz, which rotates with respect to
the inertial system XYZ. In order to transform the body-fixed co-ordinates into the
inertial co-ordinates, it is necessary to use the Euler angles w, h, and /.
As has been shown in Sect. 7.8, this transformation is based on the following
equations
x1 ¼ wsin h sin / þ hcos /

x2 ¼ wsin h cos / hsin /
x3 ¼ wcos h þ /
where, in the present case, the value of x3 is constant in time, because MO3 is equal
to zero, as has been shown in Sect. 10.2 of the present chapter. By multiplying all
10.3 Misalignment of the Thrust Vector in Inertial Co-ordinates 1217
terms of the second equation by i ¼ ð1Þ1=2 , adding these terms to the corre-
sponding terms of the first equation, and remembering that 1/i = −i, there results
x12 x1 þ ix2 ¼ ðhþ iwsin hÞðcos / i sin /Þ ¼ ðhþ iwsin hÞei/
The third equation, solved for w, yields
x3 /
w ¼
cos h
This expression of w, substituted into x12 ¼ ðhþ iwsin hÞei/ , yields
x12 ¼ ½hþ iðx3 /Þ tan hei/
The preceding equation expresses the complex angular velocity x12 ¼ x1 þ ix2
in terms of Euler’s angles between the body-fixed and the inertial co-ordinates.
It is possible to simplify the preceding equation when the spin axis of the
spacecraft does not deviate much from a fixed direction of flight. In this case, the
value of the angle h is so small that tan h h, and therefore the preceding equation
x12 ½hþ iðx3 /Þhei/
The preceding differential equation involves the two functions of time h and /.
In order to obtain a differential equation involving only one of these functions, Leon
[2] defines a complex angle of attack h12 as follows
h12 ¼ hei/
The preceding expression, differentiated with respect to time, yields
h12 ¼ hei/ i/hei/
By solving for hei/ and substituting into x12 ¼ ½hþ iðx3 /Þhei/ , there
results
x12 ¼ h12 þ ix3 h12
When x12 is a known function of time, then x12 ¼ h12 þ ix3 h12 is an ordinary
differential equation of the first order, whose unknown function is h12. As has been
shown above, this differential equation holds only when h is small.
Let us consider again the equation
x12 x1 þ ix2 ¼ ðhþ iwsin hÞðcos / i sin /Þ
After separating the real part from the imaginary part of the complex quantity
x12, this equation becomes
x12 x1 þ ix2 ¼ ðhcos / þ wsin h sin /Þ þ iðwsin h cos / hsin /Þ
The two parts, real and imaginary, of x12 are the components of hþ iw sin h
along the principal axes of inertia, respectively, 1 and 2. Therefore, multiplying the
real quantity ðhþ iwsin hÞ by e−i/ results in these two components, which are
represented by the complex quantity x1 þ ix2 .
Likewise, the real angle h multiplied by e−i/ represents the two components of
the complex quantity h12 along the principal axes of inertia, respectively, 1 and 2.
By so doing, the small angle h is represented by means of two components, as if it
were a vector. In case of an axially symmetric ðI1 ¼ I2 6¼ I3 Þ rocket not subject to
external moments ðMO1 ¼ MO2 ¼ MO3 ¼ 0Þ, the complex angular velocity x12
expressed by the following equation of Sect. 10.2
Zt
1
x12 ðtÞ ¼ x12 ð0Þeikt þ M O1 ðsÞeikðtsÞ ds
I1
0
reduces to
x12 ðtÞ ¼ x12 ð0Þeikt
where k ¼ x3 ðI1 I3 Þ=I1 . Substituting x12 ðtÞ ¼ x12 ð0Þeikt into the preceding
equation x12 ¼ h12 þ ix3 h12 leads to the following equation
h12 þ ix3 h12 ¼ x12 ð0Þeikt
Let x12 ðsÞ be the Laplace transform of x12(t). The subsidiary equation of the
equation written above is [1]:
h12 ð0Þ x12 ð0Þ

h12 ðsÞ ¼ þ
s þ ix3 ðs þ ix3 Þðs ikÞ
The inverse transform of h12 ðsÞ yields the solution of the differential equation
h12 þ ix3 h12 ¼ x12 ð0Þeikt , as follows
10.3 Misalignment of the Thrust Vector in Inertial Co-ordinates 1219
Zt
ix3 t ix3 t
h12 ðtÞ ¼ h12 ð0Þe þ x12 ð0Þe eiðk þ x3 Þs ds
0
eix3 t h i
¼ h12 ð0Þeix3 t þ ix12 ð0Þ 1 eiðx3 þ kÞt
x3 þ k
Thomson [1] computes the angle of attack hXY(t), expressed in the inertial system
XYZ, by multiplying all terms of the preceding equation by exp(ix3t), as follows
x12 ð0Þ h i
hXY ðtÞ ¼ h12 ð0Þ þ i 1 eiðx3 þ kÞt
x3 þ k

x12 ð0Þ x12 ð0Þ
¼ h12 ð0Þ þ sin½ðx3 þ kÞt þ i f1 cos½ðx3 þ kÞtg
x3 þ k x3 þ k
The real part and the imaginary part of the complex angle of attack hXY(t) are
along, respectively, the X-axis and the Y-axis. The tip of the vector hXY(t) describes
a circumference, whose radius is ½x12 ð0Þ=ðx3 þ kÞ, and whose centre C is the
complex number h12 ð0Þ þ i½x12 ð0Þ=ðx3 þ kÞ, as shown in the following figure
(redrawn from Ref. [1]).
10.4 Near-Symmetric Body of Revolution Not Subject

to Moments
As has been shown in Sect. 10.1, when the body-fixed axes x, y, and z of a spinning
rocket do not coincide with its principal axes of inertia 1, 2, and 3, then the
fundamental equation M O ¼ hO þ x hO governing the motion of a rocket
rotating about its centre of mass O, in scalar terms, is

MOx ¼ Ix xx Ixy xy Ixz xz þ Ixz xx Iyz xy þ Iz xz xy Ixy xx þ Iy xy Iyz xz xz

MOy ¼ Ixy xx þ Iy xy Iyz xz Ixz xx Iyz xy þ Iz xz xx þ Ix xx Ixy xy Ixz xz xz

MOz ¼ Ixz xx Iyz xy þ Iz xz þ Ixy xx þ Iy xy Iyz xz xx Ix xx Ixy xy Ixz xz xy
The solution of these three equations is not a simple task, even when the values
of the products of inertia Ixy, Ixz, and Iyz are small in comparison with those of the
moments of inertia Ix, Iy, and Iz.
When the rocket is a near-symmetric body of revolution, then its principal axis 3
of inertia deviates by a small angle b from its spin axis z. Since the choice of the
two transverse axes x and y is arbitrary, then it is possible to choose these axes so
that one of them (e.g. x) should be perpendicular to the plane z3 passing through O.
The position of the other two axes 1 and 2 is defined by the angle U between the
axes x and 1, as shown in the following figure.
Since the principal axes of inertia 1, 2, and 3 and the axes x, y, and z are fixed to
the rocket, then the two angles b and U are constant in time.
Let I1, I2, and I3 be the moments of inertia of the rocket about its principal axes
of inertia, respectively, 1, 2, and 3. The moments of inertia Ix, Iy, and Iz and the
products of inertia Ixy, Ixz, and Iyz of the rocket about the axes x, y, and z of the
body-fixed system can be expressed as functions of I1, I2, and I3 by using the
equations of Sect. 7.4, which are rewritten below
10.4 Near-Symmetric Body of Revolution Not Subject to Moments 1221
XX XX
Inn ¼ ‘na ‘nb Iab Ing ¼ ‘na ‘gb Iab
a b a b
where the coefficients ‘na and ‘nb are direction cosines, which are elements of the
following matrix
2 3
‘x1 ‘x2 ‘x3
4 ‘y1 ‘y2 ‘y3 5
‘z1 ‘z2 ‘z3
This matrix expresses the transformation from the principal axes of inertia 1, 2,
and 3 to the body-fixed axes x, y, and z. In the case considered here, the trans-
formation matrix depends on the angles b and U as follows
2 3 2 32 3
x cos U sin U 0 1
4 y 5 ¼ 4 sin U cos b cos Ucosb sin b 54 2 5
z sin U sin b cos Usinb cos b 3
In addition, since b is small, so that sin b b and cos b 1, then there results
2 3 2 32 3
x cos U sin U 0 1
4 y 5 4 sin U cos U b 54 2 5
z b sin U b cos U 1 3
Therefore, using the expressions of Inn and −Inη written above, the moments of
inertia Ix, Iy, and Iz and the products of inertia Ixy, Ixz, and Iyz of the rocket are
Ix ¼ ‘x1 2 I1 þ ‘x2 2 I2 þ ‘x3 2 I3 ¼ I1 cos2 U þ I2 sin2 U

Iy ¼ ‘y1 2 I1 þ ‘y2 2 I2 þ ‘y3 2 I3 ¼ I1 sin2 U þ I2 cos2 U þ I3 b2
Iz ¼ ‘z1 2 I1 þ ‘z2 2 I2 þ ‘z3 2 I3 ¼ I1 b2 sin2 U þ I2 b2 cos2 U þ I3
Ixy ¼ ‘x1 ‘y1 I1 þ ‘x2 ‘y2 I2 þ ‘x3 ‘y3 I3 ¼ ðI1 I2 Þ sin U cos U
Ixz ¼ ‘x1 ‘z1 I1 þ ‘x2 ‘z2 I2 þ ‘x3 ‘z3 I3 ¼ ðI1 I2 Þb sin U cos U
Iyz ¼ ‘y1 ‘z1 I1 þ ‘y2 ‘z2 I2 þ ‘y3 ‘z3 I3 ¼ I1 b sin2 U þ I2 b cos2 U I3 b
In case of an axially symmetric rocket (such that I1 = I2), the preceding moments
and products of inertia become
Ix ¼ I1 Ixy ¼ 0
Iy ¼ I1 þ I3 b2 Ixz ¼ 0
Iz ¼¼ I1 b2 þ I3 Iyz ¼ ðI1 I3 Þb
and the angle b becomes

Iyz
b¼
ðI3 I1 Þ
In order to determine the components of the angular velocity vector x, we first

operate a transformation of co-ordinates from xyz to 123, as follows
2 3 2 32 3
x1 cos U sin U bsinU xx
4 x2 5 4 sin U cos U bcosU 54 xy 5
x3 0 b 1 xz
This transformation, which uses the transpose RT of the rotation matrix R, yields
x1 ¼ xx cos U þ xy sin U þ xz b sin U

x2 ¼ xx sin U þ xy cos U þ xz b cos U
x3 ¼ xy b þ xz
Then we determine the complex angular velocity x12 by adding ix2 to x1, where
i ¼ ð1Þ1=2 , as follows
x12 ¼ x1 þ ix2 ¼ xx cos U þ xy sin U þ xz b sin U þ iðxx sin U þ xy cos U

þ xz b cos UÞ ¼ ½ðxx þ ixy Þ þ ibxz ðcos U i sin UÞ ¼ ½ðxx þ ixy Þ þ ibxz e
iU
ðxxy þ ibx3 ÞeiU
where xxy ¼ xx þ ixy and x3 xz . The equation written above x12

ðxxy þ ibx3 ÞeiU can be solved for xxy by multiplying all terms by eiU . This
yields
xxy x12 eiU ibx3
The following figure shows the axes x, y, and z fixed to the spinning rocket, its
principal axes of inertia 1, 2, and 3, and the inertial system XYZ. The same figure
also shows the line of nodes n between the planes xy and XY.
The x-axis of the rocket is perpendicular to the axes z and Z. The axes x, y, and
z fixed to the rocket are related to the axes X, Y, and Z of the inertial system as will
be shown below.
Starting from a position such that the axes Z and z coincide, the sequence of
rotations is as follows:
• first, a rotation w about the Z-axis, such that the new X-axis is brought to
coincide with the line of nodes n;
• then, a rotation h about the line of nodes n; and
• finally, a rotation / about the spin axis z of the rocket.
This sequence of rotations, which is based on the classical Euler angles w, h, and
/, is the 3-1-3 sequence. The position of the body-fixed axes x, y, and z with respect
to the principal axes of inertia 1, 2, and 3 is defined by means of the angles U and b.
As has been shown in Sect. 10.3, the complex angular velocity x12(t) and the
complex angle of attack h12(t), for an axially symmetric ðI1 ¼ I2 6¼ I3 Þ rocket not
subject to external moments ðMO1 ¼ MO2 ¼ MO3 ¼ 0Þ, are expressed by the fol-
lowing equations
x12 ðtÞ ¼ x12 ð0Þeik t

eix3 t h i
h12 ðtÞ ¼ h12 ð0Þeix3 t þ ix12 ð0Þ 1 eiðx3 þ kÞt
x3 þ k
where k ¼ x3 ðI1 I3 Þ=I1 . Also, in the present section, the initial value x12(0) of
the complex angular velocity x12 has been shown to be
x12 ð0Þ ¼ ½xxy ð0Þ þ ibx3 eiU

By substituting this expression into x12 ðtÞ ¼ x12 ð0Þeikt , there results
x12 ðtÞ ¼ ½xxy ð0Þ þ ibx3 eiðkt/Þ
By substituting x12 ðtÞ ¼ x12 ð0Þeikt into xxy ¼ x12 eiU ibx3 , there results
xxy ðtÞ ¼ xxy ð0Þeikt ibx3 ð1 eikt Þ
The preceding equation transforms the complex angular velocity x12(t) in order
to express it in the body-fixed co-ordinates x, y, and z.
By setting hxy ¼ hei/ , it is possible to represent the angle h as a vector whose
components are along the axes x and y of the body-fixed system xyz. Likewise, the
angle b can be represented as a vector whose unique component in the same system
is along the x-axis, because b is defined as the angle between the spin axis of the
rocket and its principal axis of inertia 3. Therefore, the components of the angles h
and b along the axes x and y are
ðh cos / þ bÞ iðh sin /Þ ¼ hei/ þ b ¼ hxy þ b
The expression written above can also be related to the principal axes of inertia
1 and 2 of the rocket. For this purpose, we multiply all terms of the preceding
expression by eiU ¼ cos U i sin U, because U is defined as the angle between the
axes x and 1. This yields
h12 ¼ ðhxy þ bÞeiU
The inverse of the preceding expression results from multiplying all terms of it
by eiU ¼ cos U þ i sin U. This yields
hxy ¼ h12 eiU b
The complex coefficient h12 which appears in the preceding equation has been
expressed as a function of time as follows
eix3 t h i
h12 ðtÞ ¼ h12 ð0Þeix3 t þ ix12 ð0Þ 1 eiðx3 þ kÞt
x3 þ k
Therefore, by multiplying all terms of this equation by eiU and then subtracting
b, there results
i
eix3 t h
hxy ðtÞ ¼ h12 ð0Þeix3 t þ ix12 ð0Þ 1 eiðx3 þ kÞt eiU b
x3 þ k
Then, by substituting h12 ð0Þ ¼ ½hxy ð0Þ þ beiU and x12 ð0Þ ¼ ½xxy ð0Þ þ ibx3 eiU
into the preceding equation, this equation becomes
xxy ð0Þeix3 t h i
hxy ðtÞ ¼ hxy ð0Þeix3 t þ i 1 eiðx3 þ kÞt
x3 þ k
ix3 t h i
ix3 t x 3 e iðx3 þ kÞt
þb e 1e 1
x3 þ k
This equation expresses the angle h, shown in the preceding figure, between the
spin axis of the rocket and the Z-axis of the inertial system XYZ, as if the angle h
were a vector, whose components are given along the axes x and y of the body-fixed
system xyz, such that x is the real axis and y is the imaginary axis.
The same angle h can also be expressed as a vector, whose components are along
the axes X and Y of the inertial system XYZ. For this purpose, all terms of the
preceding equation are multiplied by expðix3 tÞ ¼ cosðx3 tÞ þ i sinðx3 tÞ.
The resulting expression is given below.

xxy ð0Þ h i k x3 iðx3 þ kÞt
hXY ðtÞ ¼ hxy ð0Þ þ i 1 eiðx3 þ kÞt þ b eix3 t þ e
x3 þ k x3 þ k x3 þ k
This equation is compared with the corresponding equation derived in

Sect. 10.3, which is rewritten below for convenience.
x12 ð0Þ h i
hXY ðtÞ ¼ h12 ð0Þ þ i 1 eiðx3 þ kÞt
x3 þ k
This comparison shows that the vector h has, in the plane XY of the inertial
system XYZ, an additional term due to the angle b. The value of this angle differs
from zero when the principal axis of inertia 3 does not coincide with the spin axis
z of a rocket. In other words, when value of the angle b differs from zero, then a
rocket has not only a motion of precession at an angular velocity x3 + k, but also a
motion of nutation at an angular velocity x3. The resulting motion of the tip of the
spin axis of this rocket in the plane XY is shown in the following figure, which is
due to the courtesy of the University of Utah [3].
10.5 Rockets of Variable Mass
The masses of the rockets considered in the preceding sections have been assumed
to be constant in time (m` = 0). However, in most cases, the variation in time m` of
the mass m of a rocket is not negligible. The present section and the following ones
consider the effects induced in a rocket by a variation of its mass.
According to Newton’s second law of motion, the variation in time of the
momentum mv of a body is equal to the resultant f of the forces acting on it. When
the mass of a body is constant in time, then the variation of its momentum in time is
d(mv)/dt = m(dv/dt). Otherwise, it is necessary to consider the variation in time of
the momentum mv of a body whose mass is variable.
With reference to the preceding figure, we consider a body of variable mass

which moves at a velocity v with respect to an inertial system of reference. Let m be
the instantaneous mass of the body at a given time t. In the general case, the mass of
the body can vary positively (m` > 0) or negatively (m` < 0) in time.
10.5 Rockets of Variable Mass 1227
Let m + dm be the instantaneous mass of the same body at the time t + dt. Let v0
be the instantaneous velocity possessed by the increment (or decrement) of mass
dm at the time t. After the lapse of an infinitesimal interval of time dt, that is, at the
time t + dt, the mass of the body changes from m to m + dm.
At the time t, the momentum of the whole system, comprising the body and its
variation of mass, is
mv þ v0 dm
At the time t + dt, the momentum of the whole system is
ðm þ dmÞðv þ dvÞ mv þ mdv þ vdm
where the product (dm)(dv) of the two infinitesimal quantities dm and dv has been
neglected. Therefore, the rate of change of momentum in time results from sub-
tracting the first from the second of two momenta written above and dividing by
dt. This yields
ðmv þ mdv þ vdmÞ ðmv þ v0 dmÞ dv dm

¼ m þ ð v v0 Þ
dt dt dt
Since the rate of change of momentum in time is equal to the force f, then the
preceding equation shows that the force f is used not only to accelerate the mass
m of the rocket but also to change the momentum (from v0dm to vdm) of the
additional mass dm. For all rockets, the following equation holds
dm
fT ¼ u
dt
where fT is the thrust force exerted by the jet, and u = −(v − v0) is the velocity of
the exhaust gas with respect to the nozzle.
Therefore, the equation of motion for a rocket moving along a straight line can
dv
f þ fT ¼ m
dt
The force f includes the external forces acting on the rocket (e.g. the weight and
the atmospheric drag and lift), as shown in the following figure, which is due to the
courtesy of NASA [4]. As shown in this figure, the weight force is applied to the
centre of gravity, the aerodynamic forces of drag and lift are applied to the centre of
pressure, and the thrust force is applied to the centre of the exit section of the
nozzle. As a result of the distance other than zero between the centre of gravity and
the centre of pressure, a moment acts on the rocket.
If the rocket were not spinning or turning and were also subject only to forces
passing through its centre of mass, then the moment acting on it would be zero, and
its motion would be purely translational.
10.6 Damping Effect of the Exhaust Gas in a Non-spinning

Rocket of Variable Mass
When a non-spinning rocket rotates at an angular velocity x about an axis per-

pendicular to its axis of symmetry, then the exhaust gas acquires a momentum
component −m`‘x perpendicular to the axis of symmetry, where m is the instan-
taneous mass of the rocket at a given time t, m` < 0 is the rate of change of the mass
m in time, ‘ is the distance between the centre of mass O of the rocket and the exit
section of the nozzle, and ‘x is the transverse component of the velocity of the
exhaust gas at the exit section of the nozzle due to the angular velocity x of the
rocket, as shown in the following figure.
The axis of rotation of the rocket passes through its centre of mass O and is
perpendicular to the plane of the sheet. Let I = mk2 be the moment of inertia of the
10.6 Damping Effect of the Exhaust Gas in a Non-spinning … 1229
rocket about the axis of rotation specified above, where k is the radius of gyration
about the same axis.
The moment MO acting about the same axis is equal to the change in the moment
of momentum of both the rocket and the exhaust gas, which possesses a momentum
component −m`‘x perpendicular to the axis of symmetry of the rocket itself. This
fact is expressed in mathematical terms by the following equation
dðIxÞ
MO ¼ m‘2 x
dt
By executing the derivative of the product Ix (where I = mk2) with respect to

time, there results

dx dðmk 2 Þ dð k 2 Þ
MO ¼ I þx m‘2 x ¼ Ix x m ‘2 k2 m
dt dt dt
When MO = 0 and d(k2)/dt is negligible (which is the case with a

solid-propellant rocket in which the burn proceeds radially), then the preceding
differential equation can be rewritten as follows

dx ‘2 k 2 ‘2 dm
¼ dm ¼ 1
x I k 2 m
This is because I = mk2, as has been shown above. The preceding differential
equation, with the associated initial condition x = x0 and m = m0 at t = 0, has the
following solution

‘22 1
x m k
¼
x0 m0
The angular velocity x of the rocket decreases in time with respect to its initial
value x0 when the exponent (‘2/k2 – 1) is greater than zero, and increases when the
same exponent is less than zero. In most practical cases, ‘2/k2 > 1 and therefore the
gas ejected has a damping effect on the attitude of a rocket.
10.7 Euler’s Equations for Spinning Rockets of Variable

Mass
The fundamental equation
M O ¼ hO þ x hO
which governs the motion of a rigid body rotating about its centre of mass O must
be modified for a body of variable mass, as is the case with a rocket ejecting
particles of gas out of its nozzle. Considering a rocket as a system of particles, the
moment of the external forces acting on the system, taken with respect to any axis
which passes through the centre of mass of the system and is fixed in direction, is
equal to the time rate of change of the moment of momentum of the system, taken
with respect to that axis, plus the rate at which the particles which are leaving the
system are transferring moment of momentum, taken with respect to that axis, out
of the system [13, pp. 11–12], and [5]. Therefore, in order to take account of the
rate of change in time of the momentum transfer from the varying mass of
the rocket, a further term must be added to the terms (hÒ and x hO) which
are on the right-hand side of the equation written above, as follows
M O ¼ hO þ x hO + rate of change in time of the moment of momentum of the
varying mass
With reference to the following figure, let us consider the motion of an axially
symmetric rocket having a body-fixed system of reference xyz, whose origin is in
the centre of mass O of the rocket.
The particles of gas ejected by the rocket pass through a cluster of N nozzles.
The centre of each nozzle is identified by its position vector ri, which has the
following components along the axes x, y, and z
ri ¼ xi ux þ yi uy ‘uz
where i = 1, 2, …, N (N = 4 in the figure), as shown below.
Choosing the body-fixed axes x, y, and z coincident with the principal axes of
inertia 1, 2, and 3 of the rocket, the components of the moment-of-momentum
vector hO of the rocket along these axes are, respectively, hOx ¼ Ix xx ; hOy ¼ Iy xy ,
and hOz = Izxz, where Ix, Iy, and Iz are the instantaneous values of the moments of
inertia of the rocket about the principal axes of inertia, respectively, x, y, and z.
10.7 Euler’s Equations for Spinning Rockets of Variable Mass 1231
The rate of change in time of the moment-of-momentum vector of the rocket is
h0 O ¼ hO þ x hO
Taking account of the rates of change in time I`x, I`y, and I`z of the instantaneous
values of the moments of inertia of the rocket, the components of the vector h′O
along the body-fixed axes x, y, and z are respectively

Ix xx þ Ix xx Iy Iz xy xz
Iy xy þ Iy xy ðIz Ix Þxz xx

Iz xz þ Iz xz Ix Iy xx xy
As shown in the preceding figure, the cluster of nozzles is assumed to be

symmetric with respect to the longitudinal axis z of the rocket. When the angular
velocities xx, xy, and xz of the rocket are equal to zero, then the moments of
momentum of the jets are also equal to zero.
Otherwise, when xx, xy, and xz are not equal to zero, then the linear velocities
of the exit sections of the nozzles are those shown in the preceding figure. The rates
of change in time of the linear momenta relating to these linear velocities result
from multiplying the linear velocities by m` dm/dt, which is the mass flow rate of
the exhaust gas. Finally, the rates of change in time of the moments of momentum
result from multiplying the rates of change in time of the linear momenta by their
respective distances from the co-ordinate axes.
Following Ellis and McArthur [6], we consider the ith nozzle of the rocket. For
this nozzle alone, the rate of change in time of the linear momentum due to xz is
mi ri xz , and the rate of change in time of the moment of momentum about the
z-axis is mi r 2 i xz , where mi ¼ dmi =dt is less than zero, because the mass of the
rocket decreases in time. The total rate of change in time of the moment of
momentum about the z-axis for all the N nozzles is
X
N
xz mi ri2 ¼ mq2 xz
i¼1
where
PN
mi ri2
q ¼
2 i¼1
m
We consider now, for the ith nozzle of the rocket, the rate of change in time of
the linear momentum due to xx, which is mi ð‘2 þ y2i Þ1=2 xx , and the rate of
change in time of the moment of momentum about the x-axis, which is
mi ð‘2 þ y2i Þxx . The total rate of change in time of the moment of momentum
about the x-axis for all the N nozzles is
X
N 1
xx mi ‘2 þ y2i ¼ mxx ‘2 þ q2
i¼1
2
Likewise, the total rate of change in time of the moment of momentum about the
y-axis for all the N nozzles, due to xy, is
X
N 1
xy mi ‘2 þ x2i ¼ mxy ‘2 þ q2
i¼1
2
We collect the terms determined above, relating to the total rate of change in
time of the moment of momentum due to the jets about the three axes x, y, and z,
and add the components of the vector h0 O due to the rocket along the same axes.
These components of h0 O have been determined previously and are rewritten below
for convenience:

Ix xx þ Ix xx Iy Iz xy xz
Iy xy þ Iy xy ðIz Ix Þxz xx

Iz xz þ Iz xz Ix Iy xx xy
By so doing, it is possible to write the following equations

MOx ¼ Ix xx þ Ix xx Iy Iz xy xz mxx ð‘2 þ 1=2q2 Þ
MOy ¼ Iy xy þ Iy xy ðIz Ix Þxz xx mxy ð‘2 þ 1=2q2 Þ

MOz ¼ Iz xz þ Iz xz Ix Iy xx xy mxz q2
Remembering that Ix ¼ mk 2 x and substituting Ix ¼ mk2 x þ mðk2 x Þ, and similar
expressions for I`y and I`z into the three preceding equations, where kx, ky, and kz are
the radii of gyration of the rocket about the axes, respectively, x, y, and z, there
results

1 2 d kx2
MOx ¼ Ix xx Iy Iz xy xz m ‘ þ q kx m
2 2
xx
2 dt
" #
1 d ky2
MOy ¼ Iy xy ðIz Ix Þxz xx m ‘2 þ q2 ky2 m xy
2 dt

d kz2
MOz ¼ Iz xz Ix Iy xx xy m q2 kz2 m xz
dt
10.7 Euler’s Equations for Spinning Rockets of Variable Mass 1233
The preceding equations are the Euler equations for a spinning rocket of variable
mass. They contain, in comparison with the Euler equations for a rigid body,
additional terms, which are due to jet damping and the variable moments of inertia
of the rocket.
As an application, Thomson [1] considers the motion of an axially symmetric
(Ix = Iy = I) missile, not subject to external moments MOx ¼ MOy ¼ MOz ¼ 0 ,
whose initial angular velocity about the spin axis is xz0. The propellant of the
missile is assumed to burn so that the rates of change in time of the gyration radii kx,
ky, and kz are negligible. Remembering that Iz ¼ mk 2 z , the third of the three Euler
equations, modified for a spinning rocket of variable mass, can be written as follows
2
dxz q dm
¼ 1
xz kz2 m
This equation, integrated with the associated initial condition xz = xz0 and
m = m0 for t = 0, yields

q2
1
xz m kz2
¼
xz0 m0
Substituting this expression of xz/xz0 into the first and the second of the Euler
equations, multiplying all terms of the second equation by i = (−1)1/2, and adding
the result to the corresponding terms of the first equation, there results

2 q2
1
dxxy ‘ þ 1=2q2 k 2 dm 2
kz m kz2
¼ i 1 x z0 dt
xxy k2 m k2 m0
where xxy ¼ xx þ ixy . We assume the mass m of the missile to vary linearly in
time, so that
m ¼ m0 mt
Substituting this expression of m into the preceding differential equation and

integrating it over the interval [0, t], there results

2 Zt q2
1
xxy ‘ þ 1=2q2 k 2 m kz2 m kz2
ln ¼ ln i 1 xz0 1 s ds
xxy0 k2 m0 k2 m0
0
After setting 1 ðm=m0 Þs ¼ n and ½ðq2 =k 2 z Þ 1 ¼ K, the integral written on the

right-hand side of the preceding differential equation can be computed as follows

m 2 3
1
Z m0
t
1 m
t q2
m0 m0 nK þ 1 m0
m0 kz2 4 m kz2 5
nK dn ¼ ¼ 1 1
m m K þ 1 1 m q2 m0
1
Substituting this integral into the preceding differential equation, there results
2 3 2 3
‘2 þ 1 22q2 k2
=

2 2 q22
xxy m0 5 ¼ ixz0 m0 k k m z 5
ln4 1 z2 4 1
k k
z
t 1
xxy0 m0 mt m q2 k m0
Since exp[ln(x)] = x for any x > 0, the preceding equation can be written as
follows
8 2 39
‘2 þ 1 22q2 k2
=
<
m k 2 2 q22 =
xxy m0 mt k m
1 z2 41 1 t 5
k kz
0
¼ exp ixz0 z
xxy0 m0 : m q2 k m0 ;
As has been shown by Jarmolow [7], the function xxy/xxy0 oscillates in time
with decreasing amplitude, due to jet damping, and with increasing frequency.
10.8 Angle of Attack of a Rocket
The complex angular velocity xxy ¼ xx þ ixy considered in Sect. 10.7 relates to the
axes x and y, which rotate together with the rocket. In order to evaluate the angle of
attack of the rocket, it is necessary to determine the Euler angle h measured with
respect to an axis of an inertial system. For this purpose, we express the compo-
nents xx, xy, and xz of the velocity vector x acting on the rocket by means of the
Euler angles w, h, and / (see Sect. 7.8), as follows
xx ¼ wsin h sin / þ hcos /

xy ¼ wsin h cos / hsin /
xz ¼ wcos h þ /
By adding ixy to xx, there results
xxy ¼ xx þ ixy ¼ ðhþ iwsin hÞðcos / i sin /Þ ¼ ðhþ iwsin hÞei/

10.8 Angle of Attack of a Rocket 1235
By using the following equation derived in Sect. 10.7

q2
1
xz m kz2
¼
xz0 m0
and the third equation
xz ¼ wcos h þ /
of those written above, there results

q2
1
m kz2
xz0 ¼ wcos h þ /
m0
The preceding equation, solved for w`, yields

2
3
q2
1
1 6 m kz2
7
w ¼ 4xz0 /5
cos h m0
This expression of w` is substituted into
xxy ¼ ðhþ iwsin hÞei/
This substitution yields

8 2
3 9
>
< q2
1 >
=
6 m 7 kz2
xxy ¼ hþ i4xz0 /5 tan h ei/
>
: m0 >
;
8 2
3 9
>
< q2 >
=
6 m kz2 1 7
hþ i4xz0 1 /5h ei/
>
: m0 >
;
This is because the mass of the rocket has been assumed to vary linearly as a
function of time, as follows
m ¼ m0 mt
We use the transformation
hxy ¼ hei/
which is due to Leon [2]. By differentiating the preceding expression with respect to
time, there results
hxy ¼ ðh ih/Þei/
Substituting the inverse expressions
h ¼ hxy ei/
and
h ¼ hxy ei/ þ ih/
into the following equation

8 2
3 9
>
< q22 1 >
=
6 m k
7
xxy hþ i4xz0 1 /5h ei/
z
>
: m0 >
;
yields

q2
1
m kz2
xxy ¼ hxy þ ixz0 1 hxy
m0
By substituting the following expression derived in Sect. 10.7

8 2 39
‘2 þ 1 22q2 k2 <
m k 2 q2
m kz2 5=
=
xxy m0 mt k
0 kz2 4
¼ exp ixz0 z
1 2 1 1 t
xxy0 m0 : m q2 k m0 ;
into the equation written above, there results

q22 1
‘2 þ 1= 22 q2 k2
m m
kz k
hxy þ ixz0 1 hxy ¼ xxy0 1
t t
m0 m0
8 2 39
< q2

m k2
0 kz2 4 m kz2 5=
exp ixz0 z
1 2 1 1 t
: m q2 k m0 ;
10.8 Angle of Attack of a Rocket 1237
The preceding differential equation holds for a rocket whose mass varies in time.
In case of a rocket of constant mass (m` = 0), the coefficient of hxy is ixz0, which is
constant. By contrast, for a rocket of variable mass, the coefficient of hxy is

q22 1
m kz
ixz0 1 t
m0
which is a function of time. In addition, for a rocket of variable mass, the forcing
function on the right-hand side of the equation tends to decrease the component xz
of the angular velocity vector x. This effect is due to jet damping and variable mass.
The equation written above is a linear differential equation of the first order of
the following type
xþ px ¼ q
where x, its first derivative x`, p, and q are all of them functions of time. When the
functions p and q are continuous, it can be shown (see, e.g. Ref. [8]) that the
preceding differential equation has the following solution
R
lqdt þ c
x¼
l
R
where the integrating factor l ¼ exp pdt is also a function of time, and c is a
constant which depends on the initial condition.
After the angle of attack h (expressed through the transformation hxy = he−i/)
has been determined as has been shown above, it is desirable to refer this angle to
the axes X and Y of an inertial system XYZ, because the axes x and y are fixed to the
rocket and rotate with it. As has been shown by Thomson [1], in order to determine
hXY from hxy, it is necessary to multiply the latter by
8
9
>
< q22 1 >=
m kz
expðixz tÞ ¼ exp ixz0 1 t t
>
: m0 >
;
The procedure indicated above makes it possible to determine the angle of attack
h of a rocket of variable mass as a function of time.
10.9 The Motion of a Spinning Rocket with Varying

Configuration and Mass
For the bodies which have been considered hitherto, the origin of the body-fixed
system of reference coincides with the centre of mass of the body itself. The present
section considers the translational and rotational motions of a body which changes
its configuration and mass with time. For such a body, the origin of the body-fixed
system of reference does not coincide at all times with the centre of mass of the
body.
Such is the case with a rocket having motors and other moving parts or con-
taining liquid propellants. Another example is a rocket whose flexible parts are
subject to vibratory motions, and whose mass varies as a result of gas particles
ejected from the nozzles.
In order to study this problem, a rocket is considered as an aggregate of particles
whose centre of mass does not coincide with the origin of the body-fixed system of
reference. This aggregate is supposed to be confined within a given boundary, as
shown in the following figure, and has a system of co-ordinates x, y, and z, whose
origin is in O, which moves with it.
The aggregate has a variable mass, because some of its particles leave it through
its boundary. The moment of momentum of the aggregate of particles about the
moving origin O at a given time t is
X
hO ¼ ri mi R0 i
where the position vector ri of the ith particle goes from the origin O of the
body-fixed system xyz to the particle, and R′i is the velocity vector of the particle
with respect to an inertial reference system XYZ.
10.9 The Motion of a Spinning Rocket with Varying Configuration and Mass 1239
By differentiating the preceding equation with respect to time, there results

X 0
X
h0 O ¼ ri ðmi R0 i Þ þ r0 i mi R0 i
Since, as shown in the preceding figure, Ri ¼ RO þ ri , then the second term on

the right-hand side of the equation written above can be written as follows
X X 0
X
r0 i mi R0 i ¼ r0 i mi ðRO þ ri Þ ¼ R0 O mi r0 i ¼ R0 O mr0
where m is the total mass of the aggregate of particles at a given time t, and r0 is the
derivative with respect to time of the position vector r of the centre of mass of the
aggregate at the same time t in the body-fixed system xyz.
The first term on the right-hand side of the equation
X 0
X
h0 O ¼ ri ðmi R0 i Þ þ r0 i mi R0 i
is the moment about O of the force applied to the mass mi. Therefore, by using the
notation MO to designate this moment, the preceding equation can be written as
follows
X 0
MO ¼ ri ðmi R0 i Þ ¼ h0 O þ R0 O mr0
In other words, the moment MO about an arbitrary point O is equal to the rate of
change in time h0 O of the moment of momentum hO plus a term which depends on
the velocity R0 O of O and also on the velocity r0 of the centre of mass of the
aggregate with respect to O. Therefore, MO is equal to h0 O only when:
(a) R0 O is equal to the zero vector (i.e. when O is stationary); or
(b) the velocity r0 of the centre of mass of the aggregate with respect to O is equal
to the zero vector; or
(c) the two vectors R0 O and r0 are parallel.
Following Thomson [1], the moment of momentum hO of the aggregate of
particles is considered below in two different times, t and t + Dt, in order to
determine its rate of change in time h0 O , whose expression will be substituted into
the equation written above.
The following figure shows a mass mi, whose position with respect to O at a time
t is identified by the vector ri. At a later time t + Dt, the same mass occupies another
position identified by the vector ri + Dri, and is broken into two pieces, which are,
respectively, mi þ m0 i Dt and −m0 i Dt (where m0 i ¼ dmi =dt), and move one with
respect to the other at a velocity ui.
After the lapse of the interval of time Dt, the mass mi decreases, and therefore m0 i
is less than zero. At the time t + Dt, the moment of momentum of the aggregate of
particles about the moving origin O is
X
hO þ DhO ¼ ðri þ Dri Þ ðmi þ m0 i DtÞðR0 i þ DR0 i Þ
X
þ ðri þ Dri Þ ðm0 i DtÞðR0 i þ ui Þ
where ui is negative when the mass of the exhaust gas is ejected from the nozzle of
the rocket in the direction opposite to that of R0 i .
As has been shown above, at the time t, the moment of momentum of the
aggregate of particles about the moving origin O is
X
hO ¼ ri mi R0 i
The derivative of hO with respect to time is by definition
ðhO þ DhO Þ hO
h0 O ¼ lim
Dt!0 Dt
Substituting the expressions of (hO + DhO) and hO given above, and neglecting
the infinitesimals of order higher than one, there results
X X X
h0 O ¼ ri mi R00 i þ r0 i mi R0 i ri m0 i ui
X X
¼ ri mi R00 i R0 O mi r0 ri m0 i ui
This expression of h0 O is substituted into the equation
M O ¼ h0 O þ R0 O mr0
which has been derived above. This yields

X X
MO ¼ ri mi R00 i ri m0 i ui
Thomson [1] points out that the preceding equation can also be obtained from
the general equation
X 0
MO ¼ ri ðmi R0 i Þ
by noting that
ðmi R0 i Þ0 ¼ mi R00 i m0 i ui
because MO is the moment with respect to O of the external force f ¼ mv0 m0 u

applied to the rocket. The acceleration R00 i computed above relates to the inertial
system of reference XYZ. This acceleration, with respect to the body-fixed system
xyz, has the following expression
R00 i ¼ R00 O þ ri þ 2x ri þ x ri þ x ðx ri Þ
where R00 O the acceleration of the moving origin O with respect to the inertial
system XYZ, ri is the acceleration of the ith particle with respect to the body-fixed
system xyz, and 2x ri ; x ri , and x ðx ri Þ are the fictitious accelerations
(respectively, Coriolis, Euler, and centrifugal) of the P ith particle. P
This expression of R i , substituted into M O ¼ ri mi R00 i ri m0 i ui ,
00
yields
X X X
M O ¼ R00 O mr þ ri mi ri þ 2 ri ðx mi ri Þ þ ri ðx mi ri Þ
X X
þ ri ½x mi ðx ri Þ ri m0 i ui
In order to understand the physical meaning of the terms on the right-hand side
of the preceding equation, we differentiate the following equation of Sect. 7.2, with
respect to time
hO ¼ = x
where = is the matrix of inertia

2 3
Ix Ixy Ixz
= 4 Ixy Iy Iyz 5
Ixz Iyz Iz
which has been defined in Sect. 7.2. This differentiation yields
h0 O ¼ ð= xÞ0 ¼ ð= xÞþ x ð=
= xÞ ¼ = x þ = xþ x ð=
= xÞ
The three terms = x; = x, and x ð= = xÞ on the right-hand side of the

preceding equation are respectively
X X X
= x ¼ ri mi ðx ri Þ þ ri ðx mi ri Þ þ ri m0 i ðx ri Þ
X
= x ¼ ri ðx mi ri Þ
X
x ð== xÞ ¼ ri ½x mi ðx ri Þ
The preceding equation

X X X
M O ¼ R00 O mr þ ri mi ri þ 2 ri ðx mi ri Þ þ ri ðx mi ri Þ
X X
0
þ ri ½x mi ðx ri Þ ri m i ui
can be rewritten as a function of the terms expressed above, as follows

X X
M O ¼ R00 O mr þ ri mri þ 2 ri ðx mi ri Þ þ = xþ x ð=
= xÞ
X
0
ri m i ui
Since
X
= xþ x ð= = xÞ0 = x ¼ ð=
= xÞ ¼ ð= = xÞ0 ri mi ðx ri Þ
X X
ri ðx mi ri Þ ri m0 i ðx ri Þ
then the expression of MO can be written as follows

X X
M O ¼ R00 O mr þ ri mi ri þ 2 = xÞ0
ri ðx mi ri Þ þ ð=
X X X X
ri mi ðx ri Þ ri ðx mi ri Þ ri m0 i ðx ri Þ r i m0 i u i
Since the following simplification is also possible in the preceding equation

X X X
2 ri ðx mi ri Þ ri ðx mi ri Þ ¼ ri ðx mi ri Þ
then the expression of MO becomes

X X
M O ¼ R00 O mr þ ri mi ri þ ri ðx mi ri Þ þ ð= xÞ0
X X X
ri mi ðx ri Þ ri m0 i ðx ri Þ ri m0 i ui
A further simplification is possible by including the two terms

X
ri mi ðx ri Þ
X
þ ri ðx mi ri Þ
on the right-hand side of the preceding equation into a single term, by using the
following identities (see, e.g. Ref. [9]):
v u ¼ þu v
a ðb cÞ þ ðb aÞ c ¼ b ða cÞ

X X X X
ri mi ðx ri Þ þ ri ðx mi ri Þ ¼ ri ðx mi ri Þ þ ðx ri Þ mi ri
X
¼x ðri mi ri Þ
and therefore the expression of MO can be written as follows

X X
M O ¼ R00 O mr þ ri mi ri þ x = xÞ0
ðri mi ri Þ þ ð=
X X
ri m0 i ðx ri Þ ri m0 i ui
The terms on the right-hand side of the preceding equation can be interpreted as
follows. The first term R00 O mr is due to the fact that the moving origin O of the
body-fixed system xyz does not coincide with the centre of mass of the body. The
two terms = xþ x ð= = xÞ correspond to the well-known Euler equations,
0
whereas the term ð= = xÞ includes the additional term = x, which takes account
of the rate of change in time = of the matrix of inertia =, which change is due to
the change of position
P of the particles in relative motion and also to the variation of
mass. The term ri m0 i ui is the moment due to the misalignment of the thrust
vector
P with respect to the axis of symmetry of the rocket. The term
ri m0 i ðx ri Þ is the jet damping due to the motion of rotation of the rocket.
The remaining terms are due to the relative motion of the particles. The moment MO
is due to forces whose directions do not pass through the origin of the body-fixed
system. The external forces accelerate the instantaneous centre of mass and also
change the linear momentum of the ejected particles of gas, according to the fol-
lowing equation
X
f ¼ m½R00 O þ ri þ 2x ri þ x r þ x ðx rÞ m0 i ui
The first of the two terms on the right-hand side is the acceleration of the
instantaneous centre of mass, which results from the equation derived above
R00 i ¼ R00 O þ ri þ 2x ri þ x ri þ x ðx ri Þ
concerning the acceleration acting on the ith particle. In the general case, the three
scalar equations relating to the forces f are coupled with those relating to the
moments MO.
10.10 The Yo-Yo de-Spin Mechanism
In many cases, the angular velocity of a spin-stabilised satellite has to be reduced in

order to allow the proper operation of the instruments carried on board. This
de-spinning can be obtained by using either reaction jets, as described in Sect. 9.12,
or a simple device, the so-called yo-yo de-spin mechanism.
This device is capable of reducing the spin velocity of a spacecraft to zero.
De-spinning is accomplished by reeling out masses on the end of tethers (like
yo-yos) attached to the spacecraft. This slows the spin of the spacecraft, just as
spinning figure skaters increase their moment of inertia by extending their arms in
order to reduce their spin, due to the conservation of moment of momentum.
When the masses are released, the spin of the spacecraft flings them away from
the spin axis. By so doing, the spacecraft transfers some amount of its moment of
momentum to the masses, in order to reduce its spin rate to the desired value. The
masses are often released at the end of the de-spin manoeuvre.
De-spin is needed, because some final stages of rockets are spin-stabilised, and
therefore require spin rates equal to or higher than around 50 rpm to keep them
stable during firing. After the firing, the satellite carried by this stage cannot be
simply released, since this spin rate exceeds the capability of the attitude control
system of the satellite. In order to reduce the spin rate to values (about 2–5 rpm)
which are compatible with the attitude control system of the satellite, a yo-yo
de-spin device is often used. This device is illustrated in the following figure, which
is due to the courtesy of NASA-JPL [10].
10.10 The Yo-Yo de-Spin Mechanism 1245
The Pioneer III Lunar Probe was equipped with a yo-yo de-spin mechanism
[11]. According to NASA-JPL [10], two small yo-yo de-spin masses on wires were
used to reduce the spin of the New Horizons probe before its release from the
third-stage booster. After the spin rate of the probe was lowered, these masses and
their wires were released, and so are also on an escape trajectory out of the Solar
System.
The following figure shows a scheme of this device, where for simplicity only
one of the two masses m/2 each is shown.
Let /` be the spin rate of a satellite about its axis z of symmetry at the initial time
t0 in which the mass m is detached from the cylindrical surface of the satellite. After
this time, the wire unwinds, and the spin rate of the satellite gradually decreases.
The wire, after being completely unwound, is released and allowed to fly radially
away from the body of the satellite together with the mass attached to one of its
ends. By choosing properly the length of the wire, it is possible to reduce the spin
rate of a satellite to any desired value, as will be shown below.
The following figure illustrates the first phase of the process of spin reduction, in
which the wire unwinds from ‘ = 0 to its full length ‘ = ‘f. During this phase, the
wire remains tangent to the cylindrical surface of the satellite.
In the next phase, the length of the wire remains constant in time (‘ = ‘f), and the
orientation of the wire changes from tangent to perpendicular to the cylindrical
surface of the satellite, whose radius is R. The wire with its mass m is released when
it has reached this perpendicular orientation with respect to the satellite.
The preceding figure shows only one of the two wires, because the system
(comprising a satellite, two wires, and two masses m/2 each) is symmetric. We
assume that the satellite is not subject to perturbing torques and dissipation of
kinetic energy. The total kinetic energy of the system is
1 1
T ¼ Iz /2 þ m x2 þ y2
2 2
where Iz is the moment of inertia of the satellite about its spin axis (z).
In the conditions indicated above, the total kinetic energy and the total moment
of momentum of the system are, both of them, conserved.
By inspection of the preceding figure, the co-ordinates x and y of the point mass
m can be expressed as follows
x ¼ R cos h þ ‘ sin h
y ¼ R sin h ‘ cos h
These co-ordinates are differentiated with respect to time, taking account that the
length ‘ of the wire varies with time through h and / as follows
‘ ¼ Rðh /Þ
and therefore its derivative with respect to time is
‘ ¼ Rðh /Þ
This differentiation yields
x ¼ Rhsin h þ ‘sin h þ ‘hcos h ¼ ð‘ RhÞ sin h þ ‘hcos h

y ¼ Rhcos h ‘cos h þ ‘hsin h ¼ ð‘ RhÞ cos h þ ‘hsin h
Remembering that ‘` = R(h` − /`), there results
x2 þ y2 ¼ ð‘ RhÞ2 þ ð‘hÞ2 ¼ R2 /2 þ ‘2 h2
Therefore, the total kinetic energy of the system at any time t during the first
phase can be expressed as follows
1 1 1 1
T ¼ Iz /2 þ m R2 /2 þ ‘2 h2 ¼ Iz þ mR2 /2 þ m‘2 h2
2 2 2 2
On the other hand, the total kinetic energy of the system at the time t0 just before
the detachment of the mass m from the cylindrical surface of the satellite is
1
T¼ Iz þ mR2 /20
2
This is because, at t = t0, ‘ = 0. The conservation of the total kinetic energy of

the system implies that
1 1 1
Iz þ mR2 /2 þ m‘2 h2 ¼ Iz þ mR2 /20
2 2 2
After setting k2 ¼ Iz =m þ R2 , the preceding equation can be written as follows
k2 /2 þ ‘2 h2 ¼ k2 /0 2
The total moment of momentum of the system is also conserved. The equation
which expresses this fact can be written as follows
k2 /þ ‘2 h ¼ k2 /0
To show this, we take h and / as generalised co-ordinates, h` and /` as gen-

eralised momenta, and T as the Lagrangian function of the system. The equations of
motion in Lagrangian notation can be written as follows

d @T @T
¼0
dt @h @h

d @T @T
¼0
dt @/ @/
The first and the second of the equations written above yield respectively
m‘2 h mR‘h2 ¼ 0

Iz þ mR2 /þ mR‘h2 ¼ 0
By solving the first equation for mR‘h2 and substituting into the second
equation, there results

Iz þ mR2 /þ m‘2 h ¼ 0
Remembering that k2 ¼ Iz =m þ R2 , the preceding equation can be written as

follows
k2 /þ ‘2 h ¼ 0
That is,
k2 /þ ‘2 h ¼ constant ¼ k2 /0
where the quantity k2 /0 is the moment of momentum for m = 1 at t = t0. The
preceding equation expresses the conservation of the moment of momentum.
By solving the equation k2 /þ ‘2 h ¼ k2 /0 for h and substituting into the
equation k2 /2 þ ‘2 h2 ¼ k2 /0 2 , there results
1 ‘2 =k2
/ ¼ /0
1 þ ‘2 =k2
Likewise, by solving the equation k2 /þ ‘2 h ¼ k2 /0 for / and substituting

into the equation k2 /2 þ ‘2 h2 ¼ k2 /0 2 , there results
2
h ¼ /0
1 þ ‘2 =k2
where /0 is, again, the initial spin rate of the satellite.
Since / and h have been determined, then the equation ‘ ¼ Rðh /Þ can be
integrated to determine the length ‘ of the wire wound as a function of time.
This integration yields
‘ ¼ R/0 ðt t0 Þ
In other words, the wire unwinds at a constant rate ‘ ¼ R/0 .

Since the spin rate / of the satellite decreases from its initial value /0 , then a
torque Mz ¼ Iz / acts about the spin axis z of the satellite. This torque is coun-
terbalanced by the torque generated by the tension f1 in the wire, as will be shown
below. Since h and / have been determined as has been shown above, then /
can be expressed as follows
4R‘=k2
/ ¼ 2 /0
2
1 þ ‘2 =k2
On the other hand, since Mz ¼ Iz / ¼ mðk2 R2 Þ/, then the torque Mz can be

4mR‘ 1 R2 =k2
Mz ¼ /20
2 2
1 þ ‘ =k
2
The torque Mz is counterbalanced by the torque 2Rf1 =2 ¼ Rf1 , as shown in

The counterbalancing torque −Rf1 is due to the tension f1/2 acting in each of the
two wires (here the subscript 1 indicates that this tension is computed during the
first phase of the de-spin manoeuvre).
The condition Mz = −Rf1 makes it possible to determine the value of the tension
f1 as follows

4m‘ 1 R2 =k2
f1 ¼ 2 /0
2
1 þ ‘2 =k2
The tension f1 in the wire is equal to zero at t = t0, that is, when ‘ = 0. As ‘
increases from zero to its final value ‘f, the tension f1 increases, reaches its maxi-
mum value f1MAX, and then decreases.
The maximum value f1MAX of f1 can be determined by computing the first
derivative df1/d‘ of f1 with respect to ‘ and imposing the condition df1 =d‘ ¼ 0.
By so doing, the value of ‘ for which the derivative df1/d‘ is equal to zero is
found to be
1
32
‘¼ k
3
This value of ‘, substituted into the expression of f1 written above, yields

3 1 R2
f1MAX ¼ 32 mk 1 2 /20
4 k
The second phase of the de-spin manoeuvre is illustrated in the following figure.
During this phase, the wire is completely unwound, and its length is constant in
time (‘ = ‘f = constant). The wire only changes its orientation with respect to the
satellite. With reference to the preceding figure, the point mass m at one end of the
wire has the following co-ordinates
x ¼ R cos h þ ‘f cos c
y ¼ R sin h þ ‘f sin c
The derivatives of the co-ordinates x and y with respect to time are respectively
x ¼ Rhsin h ‘f csin c
y ¼ Rhcos h þ ‘f ccos c
The sum of the squares of the expressions written above is
x2 þ y2 ¼ R2 h2 þ 2R‘f hccosðh cÞ þ ‘2f c2

Since in the second phase h is equal to /, because ‘ = ‘f = constant, then the
preceding expression can be written as follows
x2 þ y2 ¼ R2 /2 þ 2R‘f /ccosðh cÞ þ ‘2f c2
Therefore, the total kinetic energy of the system during the second phase is
1 1 h i
T ¼ Iz /2 þ m R2 /2 þ 2R‘f /ccosðh cÞ þ ‘2f c2
2 2
At the end of the second phase, that is, when the wire is released from the
satellite, there results h − c = 0, and therefore the kinetic energy of the system is
1 1
1 1
T ¼ Iz /22 þ m R2 /22 þ 2R‘f /2 cþ ‘2f c2 ¼ Iz /22 þ mðR/2
2 2 2 2
þ ‘f cÞ2
where the subscript 2 indicates the second phase of the de-spin manoeuvre.
The conservation of the total kinetic energy of the system implies that
1 1 2 1
Iz /22 þ m R/2 þ ‘f c ¼ Iz þ mR2 /20
2 2 2
At the end of the second phase, the total moment of momentum of the system is
h ¼ Iz /2 þ mð‘f þ RÞðR/2 þ ‘f cÞ
The conservation of the total moment of momentum of the system implies that

Iz /2 þ m ‘f þ R R/2 þ ‘f c ¼ Iz þ mR2 /0
The equation of the total kinetic energy is solved for ‘f c=ðR/0 Þ.

This yields
" #12 12
‘f c Iz 1 /22 =/20 Iz ð1 r 2 Þ
¼ þ 1 /2 =/0 ¼ þ 1 r
R/0 mR2 mR2
where r ¼ /2 =/0 . Then, this expression of ‘f c=ðR/0 Þ is substituted into the
equation of the moment of momentum, and this equation is solved for ‘f =R þ 1.
This yields
Iz ð1 r Þ
‘f 2
þ1
þ 1 ¼ mR 12
R Iz ð1 r 2 Þ
þ1
mR2
By setting G ¼ ½Iz ð1 r Þ=½mR2 , the preceding equation can be written as

follows
2
ðG þ 1Þ= ‘f =R þ 1 ðG þ 1Þ
r¼
G
The tension f2 in the wire at the end of the second phase must counterbalance the
sum of the two parts, mR/2 2 and m‘f c2 , of the centrifugal force, and therefore
h i
f2 ¼ m R/2 2 þ ‘f c2 ¼ m/0 2 ‘f Rð/2 =/0 Þ2 =‘f þ c2 =/0 2
¼ m/0 2 ‘f ðRr 2 =‘f þ c2 =/0 2 Þ
The term c2 =/0 2 can also be written as a function of the other variables
G, r, and ‘f /R. By so doing, the preceding equation becomes
8 " #2 9
<Rr 2 G þ 1 r ‘f =R þ 1 =
f2 ¼ m/20 ‘f þ
: ‘f ‘f =R þ 1 ‘f =R ;
The tension f2 written above is the maximum tension in the wire during the
second phase of the de-spin manoeuvre, because the length of the wire is constant
(‘ = ‘f), and the angular velocity c of the wire reaches its maximum value just
before the time of release.
The equations derived above can be simplified by noting that, in many practical
cases, the value of G is about 200 or even more, so that G G þ 1. This fact makes
it possible to approximate the preceding equation as follows
2 2
ðG þ 1Þ= ‘f =R þ 1 ðG þ 1Þ ðG þ 1Þ= ‘f =R þ 1 ðG þ 1Þ
r¼
G Gþ1
By dividing the numerator and the denominator by (G + 1), there results
Gþ1
r 2 1
‘f =R þ 1
The preceding equation can also be written as follows

2
Iz 1 þ r 1= ‘f =R þ 1
2 ¼
m R þ ‘f 1r
When ‘f/R is greater than 2p, that is, when ‘f is greater than the circumference of
2
the satellite, then the term 1= ‘f þ R is negligible in comparison with 1 + r.
In this case, the preceding equation can be written as follows
Iz 1þr
2
m R þ ‘f 1r
In many practical cases, the maximum tension in the wire occurs during the first
phase of the de-spin manoeuvre, and therefore can be computed by means of the
1
equation f1MAX ¼ 3=43 =2 mk 1 R2 =k2 /2 0 . In addition, since in practice R2/k2 is
negligible in comparison with unity, the maximum tension in the wire can be
f1MAX 1:3m/2 0 k
Taking account of the equation
Gþ1
r 2 1
‘f =R þ 1
derived above, the expression of the maximum tension f2 in the wire during the
second phase of the de-spin manoeuvre can be simplified as follows
( 2 )
Rr 2 ðr þ 1Þ ‘f =R þ 1 r
f2 ¼ m/20 ‘f þ
‘f ‘f =R
The method shown above is due to Fedor [12]. The same author has considered
not only the two tip masses, but also the masses distributed along the wires, by
including one third of the mass of the wires in the tip masses.
As an example of application, proposed by Fedor [12], the S 30 Ionosphere
Probe Satellite (Explorer 8, NSSDC/COSPAR ID: 1960-014A, shown in the fol-
lowing figure, which is due to the courtesy of NASA) was launched at an initial
spin rate /0 ¼ 15p rad/s.
It is desired to reduce the spin rate of the satellite to /f ¼ 10p=3 rad/s by using
a yo-yo de-spin device. Iz = 3.39 kg m2 and R = 0.381 m are, respectively, the
moment of inertia about the axis of symmetry and the radius of the satellite.
We first compute the ratio r ¼ /f =/0 , and obtain
10
r¼ ¼ 0:222
3 15
Then we compute
Iz 1 þ r 1 þ 0:222
2 ¼ ¼ ¼ 1:57
m R þ ‘f 1 r 1 0:222
By taking ‘f = 5.23 m as the length of each wire, we compute m as follows
Iz 3:39
m¼ 2 ¼ ¼ 0:0686 kg
1:57 R þ ‘f 1:57 ð0:381 þ 5:23Þ2
A check on the value of G yields
Iz ð1 r Þ 3:39 ð1 0:222Þ
G¼ ¼ ¼ 265
mR2 0:0686 0:3812
and therefore the approximation G G þ 1 is justified. Taking account of the mass

of the wires, and using wires made of steel (whose specific weight is 7850 kgf/m3)
having a section of 1 mm2, the total mass m0 at the tips of the two wires is
m0 m0 1 7850
m0 ¼ þ ¼ 0:0686 2 5:23 106 ¼ 0:0658 kg
2 2 3 9:807
where the factor 2 is due to the presence of two wires.

The maximum tension in each of the two wires during the first phase is
12
1 m 3:39
f1MAX ¼ 1:3 /2O k ¼ 1:3 0:5 0:0686 ð15pÞ2 þ 0:3812
2 2 0:0686
¼ 697 N
The maximum tension in each of the two wires during the second phase is
( )
1 m 2 Rr2 R Rr 2
f2 ¼ /0 ‘f þ þ ð r þ 1Þ 1 þ ¼ 669 N
2 2 ‘f ‘f ‘f
Fedor [13], Cornille [14], Mentzer [15], and Bush [16] have also described
another device, called stretch yo-yo and shown in the following figure, in which a
spring is used as a portion of each wire, for the purpose of decreasing the sensitivity
of the device to errors in the initial spin rate and in the moment of inertia.
In all the methods described above, the wires are assumed to be initially wrapped
in a circle about the satellite. Therefore, the wires, during the unwinding phase of
their motion, remain tangent to this circle of winding. This holds only in the
absence of an initial coning motion of the satellite, that is, only when the
moment-of-momentum vector, the angular velocity vector, and the axis of sym-
metry of the satellite remain in a fixed direction in space. In this case, a
two-dimensional analysis of the motion is possible. By contrast, Collins [17] has
carried out a three-dimensional analysis of the motion of a yo-yo de-spin device,
which applies to the case in which the condition mentioned above is not satisfied.
References
1. W.T. Thomson, Introduction to space dynamics (Dover Publications, New York, U.S.A,
1986). ISBN 0-486-65113-4
2. H.I. Leon, Angle of attack convergence of a spinning missile descending through the
atmosphere. J. Aerosp. Sci. 25(8), 480–484 (1958), web site http://arc.aiaa.org/doi/abs/10.
2514/8.7743?journalCode=jasps
3. University of Utah, Health Sciences Library, Gun Ballistics, web site http://library.med.utah.
edu/WebPath/TUTORIAL/GUNS/GUNBLST.html
4. NASA, Glenn Research Center, web site http://www.grc.nasa.gov/WWW/k-12/VirtualAero/
BottleRocket/airplane/rktfor.html
5. J.B. Rosser, R.R. Newton, G.L. Gross, Mathematical theory of rocket flight (McGraw-Hill,
New York, 1947)
6. J.W. Ellis, C.W. McArthur, Applicability of Euler’s dynamical equations to rocket motion.
ARS J. 29(11), 863–864 (1959), web site http://arc.aiaa.org/doi/abs/10.2514/8.4931?
journalCode=arsj
7. K. Jarmolow, Dynamics of a spinning rocket with varying inertia and applied moment. J Appl
Phys. 28(3) (1957), web site http://scitation.aip.org/content/aip/journal/jap/28/3/10.1063/1.
1722736
8. P. Dawkins, Lecture notes on differential equations, Lamar University, web site http://tutorial.
math.lamar.edu/Classes/DE/Linear.aspx
9. Wikipedia, web site https://en.wikipedia.org/wiki/Triple_product
10. NASA, Jet Propulsion Laboratory, California Institute of Technology, web site http://mars.
nasa.gov/mer/mission/launch_e_sequenceDef.html
masterCatalog.do?sc=1958-008A
12. J.V. Fedor, Theory and design curves for a yo-yo de-spin mechanism for satellites, NASA
Technical Note D-708 (1961). 16 p, web site http://www.dtic.mil/cgi-bin/GetTRDoc?AD=
AD0260758
13. J.V. Fedor, Analytical theory of the stretch yo-yo for de-spin of satellites, NASA Technical
Note D-1676 (1963). 23 p, web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.gov/
19630004528.pdf
14. H.J. Cornille, A method of accurately reducing the spin rate of a rotating spacecraft, NASA
Technical Note D-1420 (1962). 13 p, web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.
gov/19620006811.pdf
15. W.R. Mentzer, Analysis of the dynamic tests of the stretch yo-yo de-spin system, NASA
Technical Note D-1902 (1963). 24 p, web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.
gov/19630011465.pdf
16. K.S. Bush, Yo-yo despin mechanisms, NASA TM X-60068, Technical report presented at the
second aerospace mechanisms symposium, San Francisco, California, USA, 4–5 May 1967,
20 p, available at the web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.gov/
19680018324.pdf
17. R.L. Collins, A three-dimensional analysis of a tangential yo-yo despin device on a rotating
body, NASA TN D-3848 (1967). 46 p, web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.
Chapter 11
Performance and Optimisation of Rockets
11.1 Performance of a Single-Stage Rocket
The present chapter deals with the optimisation of rockets for the purpose of
achieving a desired performance. A rocket is considered as a body of variable mass
accelerated by a thrust, which is obtained by gas particles ejected at high speed
from a nozzle.
Let v, m, fT, fD, and g be, respectively, the velocity vector of a rocket, its
instantaneous mass, the thrust force, the aerodynamic drag force, and the acceler-
ation of gravity acting on it. As has been shown in Sect. 4.5, in case of a rocket,
which moves along a trajectory of powered ascent in order to be inserted into an
orbit around the Earth, the equilibrium of the accelerations due to the forces
indicated above is expressed by the following equation
dv fT fD
¼ g sin c
dt m m
where c is the instantaneous value of the flight path angle of the ascent trajectory,
that is, of the angle between the velocity vector v of the rocket and the local
horizontal. Since the magnitude of the drag acceleration vector fD/m is inversely
proportional to the characteristic length of the rocket, then the term fD/m is often
negligible in comparison with fT/m in case of large rockets.
Considering the phase of vertical ascent (where c = p/2, g g0 = 9.80665 m/s2)
of a rocket, the preceding equation can be simplified as follows
dv fT
¼ g
dt m
Since a rocket loses mass during the phase of its powered ascent, then the rate of
change dm/dt of its mass in time is negative. Remembering the thrust equation of
Sect. 4.2

https://doi.org/10.1007/978-3-319-62220-0_11
1258 11 Performance and Optimisation of Rockets
dm
fT ¼ u þ ðpe p0 ÞAe
dt
where u is the constant velocity of the exhaust gas at the exit section of the nozzle,
pe is the pressure of the gas at the exit section, p0 is the external pressure, and Ae is
the area of the exit section, and also neglecting the difference of pressure pe − p0,
the thrust equation can be written as follows
dm
fT ¼ u
dt
where the minus sign is placed in front of the term on the right-hand side of the
preceding equation in order to make fT positive. Substituting this expression of fT
into the equation dv/dt = fT/m − g, there results
dm
dv ¼ u gdt
m
The preceding differential equation, integrated with the initial condition v = v0

and m = m0 at t = 0, yields
m
0
v v0 ¼ u ln gt
m
In particular, let vb0 and mb0 be, respectively, the velocity and the mass of the
rocket at the time tb0 of burnout. The preceding equation shows that the maximum
velocity attainable by a rocket in a vertical powered flight is

m0
vb0 ¼ v0 þ u ln gtb0
mb0
which is in accordance with the following equation

mprop þ mpay
v ¼ v0 þ Isp g0 ln ðgsin cÞt
mpay
derived in Sect. 4.5, where Isp = u/g0 is the specific impulse (which is also defined in
Sect. 4.5) of the rocket, g0 = 9.80665 m/s2 is the acceleration due to gravity at sea level,
g is the assumed average value of the gravitational acceleration, and mprop and mpay are
the masses of, respectively, the propellant and the payload carried by the rocket.
In case of chemical rockets, the velocity u of the exhaust gas at the exit section of the
nozzle depends on the heat energy per unit mass of the propellant and also on the
molecular weight of the exhaust gas. For a high performance in terms of vb0, it is
necessary to use a propellant of high heat energy, whose combustion generates an
exhaust gas of low molecular weight. Table 11.1, due to the courtesy of the Government
of the USA [1, p. 44], gives the specific impulse of some typical chemical propellants for
rockets.
11.1 Performance of a Single-Stage Rocket 1259
1
Table 11.1 Specific impulse of some typical chemical propellants
Propellant combinations: Isp range (sec)
Monopropellants (liquid):
Low-energy monopropellants 160–190
Hydrazine
Ethylene oxide
Hydrogen peroxide
High-energy monopropellants:
Nitromethane 190–230
Bipropellants (liquid):
Low-energy bipropellants 200–230
Perchloryl fluoride-Available fuel
Analine-Acid
JP-4-Acid
Hydrogen peroxide-JP-4
Medium-energy bipropellants 230–260
Hydrazine-Acid
Ammonia–Nitrogen tetroxide
High-energy bipropellants 250–270
Liquid oxygen-JP-4
Liquid oxygen-Alcohol
Hydrazine–Chlorine trifluoride
Very high-energy bipropellants 270–330
Liquid oxygen and fluorine-JP-4
Liquid oxygen and ozone-JP-4
Liquid oxygen–Hydrazine
Super high-energy bipropellants 300–385
Fluorine–Hydrogen
Fluorine–Ammonia
Ozone–Hydrogen
Fluorine–Diborane
Oxidiser-binder combinations (solid):
Potassium perchlorate:
Thiokol or asphalt 170–210
Ammonium perchlorate:
Thiokol 170–210
Rubber 170–210
Polyurethane 210–250
Nitropolymer 210–250
Ammonium nitrate:
Polyester 170–210
Rubber 170–210
(continued)

Propellant combinations: Isp range (sec)
Nitropolymer 210–250
Double base 170–250
Boron metal components and oxidant 200–250
Lithium metal components and oxidant 200–250
Aluminium metal components and oxidant 200–250
Magnesium metal components and oxidant 200–250
Perfluoro-type propellants 250 and above
1
Some Considerations Pertaining to Space Navigation, Aerojet-General Corp., Special Rept.
No. 1450, May 1958
Thomson [2] defines the quantity ℛ, called thrust ratio, as follows
fT
R¼
m0 g
where fT is the thrust, and m0 g is the initial weight of the rocket. Taking account of
this definition, the duration tb0 of the powered flight can be expressed as follows

m0 mb0 Isp mb0
tb0 ¼ gIsp ¼ 1
fT R m0
Substituting this expression of tb0 into the equation

m0
vb0 ¼ v0 þ u ln gtb0
mb0
derived above, there results

m0 1 mb0
vb0 ¼ v0 þ gIsp ln 1
mb0 R m0
In other words, the maximum velocity vb0 attainable by a rocket at the time tb0 of
burnout depends on the mass fraction mb0/m0, on the specific impulse Isp of the
propellant used in the rocket, and on the thrust ratio ℛ.
Setting v0 = 0, it is possible to plot vb0 versus mb0/m0, for given values of Isp and
ℛ. The following figure shows vb0 as a function of mb0/m0, where ℛ has been set
equal to 2, and Isp has three different values, which are, respectively, 350, 300, and
250 seconds.
As shown in the figure, a single-stage rocket is unfit for the purpose of placing a
satellite into an orbit around the Earth. It is necessary to find a compromise con-
cerning the values to be chosen for ℛ, mb0/m0, and Isp. In this regard, Thomson [2]
notes that a large value of ℛ implies a heavy structure, and a large value of Isp
implies a large value of mb0/m0. In practice, it is difficult to achieve a value less than 0.1
for mb0/m0 and a value greater than 350 seconds for Isp in case of chemical propellants.
These facts indicate the necessity of using multi-stage rockets for Earth satellites and
space missions. These rockets will be considered in the following paragraphs.
In order to determine the height hb0 reached by a rocket at the time tb0 of burnout,
it is necessary to integrate the equation, derived above, expressing its velocity
m
0
v v0 ¼ u ln gt
m
During the powered ascent, we assume the variation of g with altitude to be

negligible (g g0) and the variation m′ of m with time to be known. For this
purpose, it is reasonable to assume a constant rate of propellant consumption, which
corresponds to a constant thrust. In mathematical terms, we substitute
dm m0 mb0
dt ¼ ; m0 ¼ ¼ constant
m0 tb0
into the equation of the velocity written above and then integrate from 0 to tb0. To
this end, we first evaluate the following integral
Ztb0 Zmb0
1
lnðmÞdt ¼ 0 lnðmÞdm
m
0 m0
An integration by parts yields

Z Z
1
lnðmÞdm ¼ m lnðmÞ m dm ¼ m lnðmÞ m
m
Therefore, remembering that dt = (1/m′)dm, with m′ = constant, there results
Ztb0
1h imb0 m m
0 b0 m0 mb0
lnðmÞdt ¼ 0
m lnðmÞ m ¼ 0 lnðm0 Þ þ 0 lnðmb0 Þ
m m0 m0 m m
0
By multiplying the terms on the two sides of the preceding equation by u and
substituting m0 ¼ ðm0 mb0 Þ=tb0 into it, there results
Ztb0
m0 mb0
u lnðmÞdt ¼ utb0 1 þ lnðm0 Þ lnðmb0 Þ
m0 mb0 m0 mb0
0
Therefore,
Ztb0 Ztb0
m0
u ln dt ¼ utb0 lnðm0 Þ u lnðmÞdt
m
0 0

m0 mb0
¼ utb0 lnðm0 Þ utb0 1 þ lnðm0 Þ lnðmb0 Þ
m0 mb0 m0 mb0
Simplifying
R the equation
R written above and taking account of the other two
integrals v0 dt = v0t and gtdt ½ g0t2 (because g g0) evaluated from 0 to tb0,
the height hb0 reached by the rocket at the time tb0 of burnout is

1 m0 1 2
hb0 ¼ utb0 1 ln þ v0 tb0 g0 tb0
ðm0 =mb0 Þ 1 mb0 2
2
Isp mb0 1 m0 Isp mb0
¼ g0 1 1 ln þ v0 1
R m0 ðm0 =mb0 Þ 1 mb0 R m0
2 2
1 Isp mb0
g0 2 1
2 R m0
After the time tb0 of burnout, the rocket moves along a free-fall trajectory,
because it is subject only to the force of gravity of the Earth. From this time
onwards, the height reached by the rocket over the surface of the Earth is so great as
to make it necessary to take account of the variation of g. In order to compute g as a
function of altitude, we approximate the Earth to a sphere, whose radius R is taken
equal to the mean radius of the Earth. At the sea level, the acceleration of gravity is
lE
g0 ¼
R2
where lE = 398600.4415 km3/s2 is the gravitational parameter of the Earth.

At an altitude h above the sea level, the acceleration of gravity is
lE lE
g¼ ¼
r 2
ð R þ hÞ 2
where r = R + h is the distance of the rocket from the centre of mass of the Earth.
The ratio of g to g0 is
2
g R
¼
g0 Rþh
Therefore, g can be expressed approximately as a function of h as follows

2
R
g ¼ g0
Rþh
Since the gravitational field is conservative, then the kinetic energy ½mb0v2b0
possessed by the rocket at the time tb0 of burnout can be equated to the work done
by the force of gravity. This force is
2 2
R R
f ¼ mb0 g ¼ mb0 g0 ¼ mb0 g0
Rþh r
because R + h = r. The work done by the force of gravity from the initial radius rb0
of burnout to the final radius rb0 + hc is
Z þ hc
rb0 Zþ hc
rb0 2
R
f dr ¼ g0 mb0 dr
r
rb0 rb0
By equating the kinetic energy to the work, there results
Z þ hc
rb0 2
1 R
mb0 v2b0 ¼ g0 mb0 dr
2 r
rb0
Since
Z þ hc
rb0 2 rb0 þ hc
R 1 1 1
dr ¼ R2 ¼ R2
r r rb0 rb0 rb0 þ hc
rb0
then, substituting this integral into the equation of the kinetic energy, there results

1 2 1 1
v ¼ g0 R2

2 b0 rb0 rb0 þ hc
The preceding equation can be solved for hc, taking account that rb0 = R + hb0.
This yields
v2b0 ðR þ hb0 Þ2
hc ¼ 2
2g0 R2 vb0 =ð2g0 Þ ðR þ hb0 Þ
The total height hT, reached by the rocket after travelling along the two segments
(powered and free-fall) of its trajectory, results from
hT ¼ hb0 þ hc
where hb0 and hc can be computed by means of the equations given above.
These equations indicate that the performance of a single-stage rocket depends
on the specific impulse Isp, on the thrust ratio ℛ, and on the mass ratio l = m0/mb0
of that rocket. The values to be given to these quantities must be chosen by the
designer. The effect resulting from a variation of these quantities on vb0 or on hb0
can be evaluated by writing the two equations expressing, respectively, vb0 and hb0
as functions of Isp, ℛ, and l, as follows

vb0 ¼ f 1 Isp ; R; l

hb0 ¼ f 2 Isp ; R; l
By differentiating the functions f1 f1(Isp, ℛ, l) and f2 f2(Isp, ℛ, l) with

respect to the three variables Isp, ℛ, and l, there results
@f1 @f1 @f1

dvb0 ¼ dIsp þ dR þ dl
@Isp @R @l
@f2 @f2 @f2
dhb0 ¼ dIsp þ dR þ dl
@Isp @R @l
A necessary condition for the optimum velocity of the rocket at the time tb0 of
burnout is dvb0 = 0. Likewise, a necessary condition for the optimum height of the
rocket at the same time is dhb0 = 0. These two conditions define the constraint to be
imposed to the variables Isp, ℛ, and l.
11.2 Multi-stage Rockets
As has been shown in Sect. 11.1, a single-stage rocket is not adequate for the
purpose of placing an artificial satellite into an orbit around the Earth.
For a typical chemical single-stage rocket having a specific impulse Isp = 300 s,
a mass ratio mb0/m0 = 0.1, and a thrust ratio ℛ = 2, the velocity at the time of
burnout, in case of v0 = 0, results from the following equation derived in Sect. 11.1

m0 1 mb0
vb0 ¼ v0 þ gIsp ln 1
mb0 R m0
where g g0 = 9.80665 m/s2, which gives vb0 = 5.45 km/s. On the other hand, the
velocity of a satellite in a circular orbit of low altitude (h = 200 km) above the
surface of the Earth is
l 12 12
398600
v¼ E
¼ ¼ 7:78 km/s
r 6378 þ 200
Therefore, in order to reach this goal, it is necessary to use either a single-stage

rocket having a higher specific impulse (above 400 s, as is the case with the liquid
hydrogen–liquid oxygen engines of the Space Shuttle) or a multi-stage rocket.
A multi-stage rocket is a rocket having two or more stages, placed one on the top
of another, whose engines are fired successively. In a multi-stage rocket, the
velocity vb0 of the first stage at burnout time is added to the initial velocity
impressed by the motor of the second stage, and so on with the further stages. This
offers the advantage of eliminating dead weight, because the first stage with its
engines can be dropped off at burnout, when its fuel tank is empty. This decreases
the mass which the remaining stages of the rocket must carry. In addition, since the
huge thrust required to lift the full-weight rocket from its launch pad is no more
necessary, then the thrusts provided by the engines of the further stages can be
lower than the thrust provided by the engine of the first. Such has been the case with
the three-stage Saturn V launch vehicle, which took astronauts to the Moon. This
rocket and its performance are shown in the following figure, which is due to the
courtesy of the Smithsonian National Air and Space Museum [3].
The process of discarding a portion of the rocket after the powered ascent is
called staging. Generally speaking, there are two types of rocket staging.
The first type, called serial (or vertical) staging, is shown in the following figure,
due to the courtesy of NASA [4], which illustrates a two-stage rocket.
In this type of staging, there are one or more upper stages mounted serially on
the top of the first stage, which is the largest of all. By the way, a common method
for obtaining a first stage of high thrust is to cluster several rockets together, which
are fed by a single set of propellant tanks. The first stage is ignited at the time t = 0
of lift-off and burns through the time tb0 of powered ascent, that is, until its pro-
pellant is exhausted. After the outboard engines receive a cut-off signal, the
retro-rockets of the first stage are ignited. Then, the first stage separates from the
second, and the engine of the second stage is ignited. The same procedure is
repeated in case of the presence of a third stage.
11.2 Multi-stage Rockets 1267
The following figure, due to the courtesy of NASA [5], shows the ignition of the
ejection charge for the separation of two successive stages in case of a
solid-propellant rocket.
The same figure also shows a “regular” solid-propellant rocket motor equipped
with a delay charge and an ejection charge for separation, and a “booster”
solid-propellant rocket motor without such charges. Further information on stage
separation mechanisms used for rockets can be found in Ref. [6].
The second type, called parallel staging, is shown in the following figure, also
due to the courtesy of NASA [4], which illustrates two small boosters, called
strap-ons, placed at the sides of a longer-burning central sustainer rocket.
At the moment of lift-off, all the engines of the rocket (comprising the sustainer
and the strap-ons) are ignited simultaneously. During the powered ascent, when the
propellants in the two strap-ons are extinguished, these small rockets are discarded.
The propellant in the sustainer continues to burn until the payload carried on the top
is inserted into the desired orbit. Parallel staging is used on the Space Shuttle, which
comprises the orbiter, the main tank, and two solid-propellant boosters. These
boosters, after being dropped off, can be retrieved from the ocean, refilled with
propellant, and used again on the Shuttle.
Some rockets, for example the Titan III and the Delta II, use both serial and
parallel staging.
Care must be taken in order for the thrust to be balanced around the central line
of the rocket. In addition, as shown above, the clustered stages of the sustainer and
the strap-ons must be ignited simultaneously, in order for the rocket to attain the
highest thrust at the moment of lift-off.
11.3 Optimum Staging for Multi-stage Rockets
In this paragraph, we consider an N-stage rocket of the serial type and assign a specific
impulse (Isp)i and a structural ratio ei = [ms/(ms + mprop)]i, where i = 1, 2, …, N, to each
of its stages. We want to determine the N-stage rocket of minimum mass which carries
a given payload of mass mpay to a given velocity of burnout vb0. This problem can be
solved by using the classical method of the Lagrange multipliers, whose fundamental
concepts are shown below.
Following Curtis [7], let us consider a surface whose equation is z = f(x, y),
where f f(x, y) is a function of the two variables x and y defined on a domain R of
x and y. A given point P of the surface z = f(x, y) is said to be stationary, when it is
an extremal point, that is, either a local maximum or a local minimum of the
surface. In order for P to be stationary, the differential df of f(x, y) must be equal to
zero, that is,
@f @f
df ¼ dx þ dy ¼ 0
@x @y
where dx and dy are independent and not necessarily equal to zero. The preceding
condition implies
@f @f
¼ ¼0
@x @y
Let g g(x, y) = 0 be the equation of a curve belonging to the plane xy. We

search those points of the curve g = 0, which are stationary for z = f(x, y), and
exclude any other points of the plane xy. This is a constraint which the stationary
points of the surface z = f(x, y) must satisfy, because the two variables x and y are
11.3 Optimum Staging for Multi-stage Rockets 1269
not mutually independent: for any given value of the variable x, the variable y can
assume only those values for which the condition g(x, y) = 0 is satisfied. Since g = 0,
then dg = 0, that is,
@g @g
dg ¼ dx þ dy ¼ 0
@x @y
Since the two conditions df = 0 and dg = 0 must be satisfied at a given point, then
dy @f =@x @g=@x
¼ ¼
dx @f =@y @g=@y
After setting dy/dx = k, where k is called the Lagrange multiplier, the preceding
equation may be written as follows
@f =@x @f =@y
¼ ¼ k
@g=@x @g=@y
which in turn may be written as follows
@f @g
þk ¼0
@x @x
@f @g
þk ¼0
@y @y
These differential equations, together with the equation g(x, y) = 0 which

expresses the constraint, are the three conditions to be satisfied in order for the
function
h hðx; y; kÞ ¼ f ðx; yÞ þ kgðx; yÞ
to have a stationary point. In other terms, these three conditions are
@h @f @g
¼ þk ¼0
@x @x @x
@h @f @g
¼ þk ¼0
@y @y @y
@h
¼g¼0
@k
An extremal point P is a local maximum or a local minimum for the function

h defined above depending on whether the second differential d2h of h
@2h 2 @2h 2 @2h

d2 h ¼ dx þ dx þ 2 dx dy
@x2 @y2 @x@y
be, respectively, less than zero or greater than zero for all dx and dy satisfying the
constraint condition
gðx; yÞ ¼ 0
The method shown above is not limited to a function f of two variables x and y. It
can also be applied to the case of a function f of any number N of variables. In other
words, the general problem of optimisation can be solved as follows. According to
Coleman [8], let
f f ðx1 ; x2 ; . . .; xN Þ
be a function of N variables x1, x2, …, xN defined in a domain R. A necessary

condition for the function f to have a stationary value in R is df = 0, that is,
@f @f @f
df ¼ dx1 þ dx2 þ þ dxN ¼ 0
@x1 @x2 @xN
for all permissible values of the differentials dx1, dx2, …, dxN. This holds when the
N variables x1, x2, …, xN are independent. When these variables are not independent
but are related by a condition of the form g g(x1, x2, …, xN) = 0, then the
conditions for f to have a stationary value are
@f @g
þk ¼0
@x1 @x1
@f @g
þk ¼0
@x2 @x2
..
.
@f @g
þk ¼0
@xN @xN
g¼0
Coming back to the problem of the optimum staging for a rocket, let mi be the
so-called step mass of the ith stage of an N-stage rocket, where i = 1, 2, …, N. This
step mass is defined as follows
mi ¼ ðms þ mprop Þi
where (ms)i and (mprop)i are, respectively, the structural mass and the propellant
mass of the ith stage.
By the way, the structural mass of a stage is the mass of the structure, the
engines, the tanks of propellant, the control systems, etc., of that stage.
Let ei be the structural ratio of the ith stage, defined as follows

ms ms
i ¼ ¼
ms þ mprop i m0 mpay i
where (m0)i and (mpay)i are, respectively, the take-off mass and the payload mass of
the ith stage. We also remember the definition of mass ratio li, which has been
given in Sect. 11.1, that is

m0
li ¼
mb0 i
where (mb0)i is the burnout mass of the ith stage. Assuming that the propellant of
any stage is completely consumed at burnout, there results

m0 ms þ mprop þ mpay
li ¼ ¼
mb0 i ms þ mpay i
Remembering the definition ei = [ms/(ms + mprop)]i of structural ratio given

above, it is easy to verify that, for any stage i, the following equality holds

ðm s Þi ¼ mprop
1 i
Likewise, the structural mass (ms)i of the ith stage can be expressed as a function
of its step mass mi and its structural ratio ei, as follows

ðms Þi ¼ e ms þ mprop i ¼ ei mi
The total mass mT of an N-stage rocket without the payload results from sum-
ming the step masses of its stages, as follows
X
N
mT ¼ mi
i¼1
The total mass of the rocket with the payload at lift-off is
m0 ¼ mT þ mpay
We want to use the method of the Lagrange multipliers to determine the N-stage
rocket of minimum mass m0 at lift-off which carries a given payload of mass mpay to
a given velocity of burnout vb0. In other words, we want to determine the minimum
value of m0/mpay. We consider first a two-stage rocket (N = 2) and then a rocket
with an arbitrary number N of stages.
For N = 2, the total mass of the rocket with the payload at lift-off is
m0 ¼ m1 þ m2 þ mpay
and therefore
m0 m1 þ m2 þ mpay m1 þ m2 þ mpay m2 þ mpay
¼ ¼
mpay mpay m2 þ mpay mpay
The mass ratio of the first stage is

m0 m1 þ m2 þ mpay m1 þ m2 þ mpay
l1 ¼ ¼ ¼
mb0 1 ðms Þ1 þ m2 þ mpay 1 m1 þ m2 þ mpay
this is because (ms)i = [e(ms + mprop)]i = ei mi, and therefore, (ms)1 = e1 m1.
Likewise, the mass ratio of the second stage is

m0 m2 þ mpay m2 þ mpay
l2 ¼ ¼ ¼
mb0 2 ðms Þ2 þ mpay 2 m2 þ mpay
The equation expressing l2 can be solved for the step mass m2. This yields
l2 1
m2 ¼ mpay
1 l 2 2
Likewise, the equation expressing l1 can be solved for the step mass m1. This
yields
l1 1

m1 ¼ m2 þ mpay
1 l1 1
Let us consider the following equality

1
m1 þ m2 þ mpay 1 1 m1 þ m2 þ mpay 1 m1 þ m2 þ mpay
¼
m2 þ mpay 1 1 m2 þ mpay þ ð1 m1 1 m1 Þ 1
1 m1 þ m2 þ mpay

1

m1 þ m2 þ mpay ð1 1 Þ m1 þ m2 þ mpay 1 m1 þ m2 þ mpay
¼

m2 þ mpay 1 m1 þ m2 þ mpay 1 m1 þ m2 þ mpay 1
1 m1 þ m2 þ mpay
m1 þ m2 þ mpay
ð1 1 Þ
1 m1 þ m2 þ mpay
¼ m þm þm m1 þ m2 þ mpay
1 1 2 pay
1
1 m1 þ m2 þ mpay 1 m1 þ m2 þ mpay
Remembering that

m0 m1 þ m2 þ mpay m1 þ m2 þ mpay
l1 ¼ ¼ ¼
mb0 1 ðms Þ1 þ m2 þ mpay 1 m1 þ m2 þ mpay
there results
m1 þ m2 þ mpay ð1 1 Þl1
¼
m2 þ mpay 1 1 l 1
m2 þ mpay ð1 2 Þl2
¼
mpay 1 2 l2
Therefore, the equation derived above
m0 m1 þ m2 þ mpay m1 þ m2 þ mpay m2 þ mpay

¼ ¼
mpay mpay m2 þ mpay mpay
m0 ð1 1 Þl1 ð1 2 Þl2
¼
mpay 1 1 l1 1 2 l2
Taking the natural logarithms of the terms on both sides of the preceding
equation, there results

m0 ð1 1 Þl1 ð1 2 Þl2
ln ¼ ln þ ln ¼ lnð1 1 Þ þ lnðl1 Þ
mpay 1 1 l1 1 2 l2
lnð1 1 l1 Þ þ lnð1 2 Þ þ lnðl2 Þ lnð1 2 l2 Þ
When mpay is kept constant, then ln(m0/mpay) is a function of m0, which function
increases monotonically with m0. This is because

d m0 1
ln ¼ [0
dm0 mpay m0
Therefore, ln(m0/mpay) has a stationary value when m0 has a stationary value.

Remembering the following expression of Sect. 11.1

m0 1 mb0
vb0 ¼ v0 þ g0 Isp ln 1
mb0 R m0
and neglecting the loss (1 − mb0/m0)/ℛ, the burnout velocity for a N-stage rocket
can be expressed as follows
X
N
vb0 ¼ ui lnðli Þ
i¼1
where, for the ith stage, ui = g0 (Isp)i is the constant velocity of the exhaust gas at the
exit section of the nozzle, and li = (m0/mb0)i is the mass ratio.
In the present case of a two-stage rocket, the burnout velocity is
vb0 ¼ u1 lnðl1 Þ þ u2 lnðl2 Þ
Therefore, given the value of vb0, the equation expressing the constraint is
vb0 u1 lnðl1 Þ u2 lnðl2 Þ ¼ 0
The two equations derived above

ln m0 =mpay ¼ lnð1 e1 Þ þ lnðl1 Þ lnð1 e1 l1 Þ þ lnð1 e2 Þ þ lnðl2 Þ lnð1 e2 l2 Þ
vb0 u1 lnðl1 Þ u2 lnðl2 Þ ¼ 0
specify the functions
f f ðl1 ; l2 Þ ¼ lnð1 e1 Þ þ lnðl1 Þ lnð1 e1 l1 Þ þ lnð1 e2 Þ þ lnðl2 Þ lnð1 e2 l2 Þ

g gðl1 ; l2 Þ ¼ vb0 u1 lnðl1 Þ u2 lnðl2 Þ
The three equations to be solved are then
@f @g 1 1 u1
þk ¼ þ k ¼0
@l1 @l1 l1 1 1 l1 l1
@f @g 1 2 u2
þk ¼ þ k ¼0
@l2 @l2 l2 1 2 l2 l2
g ¼ vb0 u1 lnðl1 Þ u2 lnðl2 Þ ¼ 0
The first two of these equations, solved for, respectively, l1 and l2, yield
k u1 1 k u2 1
l1 ¼ l2 ¼
ku1 1 ku2 2
Substituting these expressions of l1 and l2 into the third equation yields

ku1 1 ku2 1
vb0 ¼ u1 ln þ u2 ln
ku1 1 ku2 2
which must be solved iteratively for k. This done, the value of k found in this
manner is substituted into the preceding equations
ku1 1 ku2 1
l1 ¼ l2 ¼
ku1 1 ku2 2
in order to determine the values of the mass ratios l1 and l2 for, respectively, the
first stage and the second stage of the rocket.
These values, in turn, are substituted into the preceding equations
l2 1 l1 1

m2 ¼ mpay m1 ¼ m2 þ mpay
1 l2 2 1 l1 1
in order to determine the values of the step masses m1 and m2 for the two stages.
The two structural ratios e1 = [ms/(ms + mprop)]1 and e2 = [ms/(ms + mprop)]2, the
two velocities u1 and u2 of the exhaust gases, and the mass mpay of the payload have
values chosen by the designer.
The method shown above for a two-stage rocket can be applied to a rocket
having an arbitrary number N of stages, as will be shown below.
In case of a rocket having N stages, the two functions defined above
f f ðl1 ; l2 Þ ¼ lnð1 e1 Þ þ lnðl1 Þ lnð1 e1 l1 Þ þ lnð1 e2 Þ þ lnðl2 Þ lnð1 e2 l2 Þ

g gðl1 ; l2 Þ ¼ vb0 u1 lnðl1 Þ u2 lnðl2 Þ
depend on the N variables l1, l2, …, lN, as follows
X
N
f ¼ lnð1 i Þ þ lnðli Þ lnð1 i li Þ
i¼1
X
N
g ¼ vb0 ui lnðli Þ ¼ 0
i¼1
The data known in the process of optimisation are the velocity vb0 at burnout, the
mass mpay of the payload, the structural ratios e1, e2, …, eN, and the velocities u1, u2,
…, uN of the exhaust gases for the N stages of the rocket.
The velocities u1, u2, …, uN depend on the specific impulses (Isp)1, (Isp)2, …,
(Isp)N through ui = (Isp)i g0 (i = 1, 2, …, N).
The first step of the computation is the determination of the Lagrange multiplier
k, by solving iteratively the general equation
X
N
kui 1
vb0 ¼ ui ln
i¼1
kui i
X
N X
N X
N
vb0 ¼ ui lnðkui 1Þ lnðkÞ ui ui lnðui i Þ
i¼1 i¼1 i¼1
After solving the preceding equation for k, the optimum values of the mass ratios
l1, l2, …, lN for the N stages result from the following N equations
kui 1
li ¼ ði ¼ 1; 2; . . .; N Þ
kui i
Since l1 = (m0/mb0)1, l2 = (m0/mb0)2, …, lN = (m0/mb0)N, then their values must

be, all of them, greater than unity.
This done, the optimum values of the step masses m1, m2, …, mN are computed
in descending order (from mN to m1), as follows
l 1
mN ¼ 1 N l mpay
N N
l 1

mN1 ¼ 1 N1lN1 N1 mN þ m pay
lN2 1

mN2 ¼ 1 l mN þ mN1 þ mpay
N2 N2
..
.
l1 1

m1 ¼ 1 l1 1 mN þ mN1 þ þ m2 þ mpay
Since the N step masses m1, m2, …, mN are known, the corresponding
N structural masses (ms)1, (ms)2, …, (ms)N result from
ðms Þi ¼ ei mi ði ¼ 1; 2; . . .; N Þ
and the N masses of propellant (mprop)1, (mprop)2, …, (mprop)N result from

mprop i ¼ mi ðms Þi ði ¼ 1; 2; . . .; N Þ
The function
h hðl1 ; l2 ; . . .; lN ; kÞ ¼ f ðl1 ; l2 ; . . .; lN Þ þ kgðl1 ; l2 ; . . .; lN Þ

where
X
N
f ðl1 ; l2 ; . . .; lN Þ ¼ lnð1 i Þ þ lnðli Þ lnð1 i li Þ
i¼1
X
N
gðl1 ; l2 ; . . .; lN Þ ¼ vb0 ui lnðli Þ ¼ 0
i¼1
is such that
@2h
¼0 ði; j ¼ 1; 2; . . .; N; i 6¼ jÞ
@li @lj
Therefore, the second differential of h, that is,
XN X N
@2h
d2 h ¼ dl dl
i¼0 j¼0
@li @lj i j
reduces to
X
N
@2h
d2 h ¼ dl2i
i¼0
@l2i
In addition, since
@ 2 h kui ð1 i li Þ2 þ 2i li 1
¼
@l2i l2i ð1 i li Þ2
and the denominator l2i (1 − liei)2 is greater than zero, then d2h is greater than zero
when
kui ð1 i li Þ2 þ 2i li 1 [ 0 ði ¼ 1; 2; . . .; N Þ
The value of k which satisfies the condition dh = 0 and also the condition written
above yields the minimum value of the mass m0 at lift-off. The method shown
above is due to Hall and Zambelli [9].
As an example of application, we search the rocket having the minimum value of
mass at lift-off, which is required to raise a payload mpay = 7000 kg to a final
velocity vb0 = 10 km/s, and which has three stages (N = 3) such that:

sp Þ1 ¼ 400 s
ðI e1 ¼ 0:10
I ¼ 350 s e2 ¼ 0:15

sp 2
Isp 3 ¼ 300 s e3 ¼ 0:20
The correspondent velocities of the exhaust gases are
u1 ¼ 400 9:80665 ¼ 3922:66 m=s 3:923 km=s

u2 ¼ 350 9:80665 ¼ 3432:33 m=s 3:342 km=s
u3 ¼ 300 9:80665 ¼ 2942:00 m=s ¼ 2:942 km=s
After substituting these data into the following equation
X
3 X
3 X
3
vb0 ¼ ui lnðkui 1Þ lnðkÞ ui ui lnðui i Þ
i¼1 i¼1 i¼1
there results
10 ¼ 3:923 lnð3:923k 1Þ þ 3:432 lnð3:432k 1Þ þ 2:942 lnð2:942k 1Þ

ð3:923 þ 3:432 þ 2:942Þ lnðkÞ 3:923 lnð3:923 0:10Þ
3:432 lnð3:432 0:15Þ 2:942 lnð2:942 0:20Þ
After simplification, the function g g(k) defined by the preceding equation is
g ¼ 2:490 3:923 lnð3:923k 1Þ 3:432 lnð3:432k 1Þ

2:942 lnð2:942k 1Þ þ 10:30 lnðkÞ
We search a value of k which makes the function g equal to zero. This value is
sought within the interval k2 k k1, where k2 = 0.4 and k1 = 0.5.
At the two ends of this interval, the values of g are
g2 gð0:4Þ ¼ 3:747
g1 gð0:5Þ ¼ 1:134
Since the condition g2g1 < 0 is satisfied, there is at least one point k within the
interval 0.4 k 0.5 such that g(k) = 0.
We choose arbitrarily another point k0 = 0.45 within the interval specified above
and compute
g0 gð0:45Þ ¼ 0:7179
Now, we set
h1 ¼ k1 k0 ; ¼ 0:50 0:45 ¼ 0:05

h2 ¼ k0 k2 ¼ 0:45 0:40 ¼ 0:05
c ¼ h2 =h1 ¼ 0:05=0:05 ¼ 1

a ¼ ½cg1 g0 ð1 þ cÞ þ g2 = ch21 ð1 þ cÞ ¼ ch21 ð1 þ cÞ

¼ ½1 ð1:134Þ 0:7179 ð1 þ 1Þ þ 3:747= 1 0:052 ð1 þ 1Þ ¼ 235:4

b ¼ g1 g0 Ah21 =h1 ¼ ð1:134 0:7179 235:4 0:052 Þ=0:05 ¼ 48:81
c ¼ g0 ¼ 0:7179
of the interpolating parabola g(k) = a(k − k0)2 + b(k − k0) + c.

Then, k is estimated as follows
h
1=2 i
k ¼ k0 2c= b b2 4ac ¼ 0:45
h i
2 0:7179= 48:81 ð48:812 4 235:4 0:7179Þ1=2 ¼ 0:4659
to minus because b has a negative value.
The correspondent value of g is g g(0.4659) = 0.04352.
For the next iteration, we reset the subscripts as follows
k2 ¼ 0:4500 g2 ¼ 0:7179
k0 ¼ 0:4659 g0 ¼ 0:04352
k1 ¼ 0:5000 g1 ¼ 1:134
and compute the corresponding values of h1, h2, and c, as follows
h1 ¼ k1 k0 ¼ 0:5000 0:4659 ¼ 0:0341

h2 ¼ k0 k2 ¼ 0:4659 0:4500 ¼ 0:0159
c ¼ h2 =h1 ¼ 0:0159=0:0341 ¼ 0:4663
Then, we compute the coefficients

a ¼ ½cg1 g0 ð1 þ cÞ þ g2 = ch21 ð1 þ cÞ
¼ ½0:4663 ð1:134Þ 0:04352 ð1 þ 0:4663Þ þ 0:7179

= 0:4663 0:03412 ð1 þ 0:4663Þ ¼ 157:6

b ¼ g1 g0 Ah21 =h1 ¼ ð1:134 0:04352 157:6 0:03412 Þ=0:0341 ¼ 39:91
c ¼ g0 ¼ 0:04352
of the interpolating parabola and estimate k as follows

h
1=2 i
k ¼ k0 2c= b b2 4ac ¼ 0:4659 2
0:04352=½39:91 ð39:912 4 157:6 0:04352Þ1=2 ¼ 0:4670
The correspondent value of g is g g(0.4670) = 0.0002737.

For the next iteration, we reset the subscripts as follows
k2 ¼ 0:4659 g2 ¼ 0:04352
k0 ¼ 0:4670 g0 ¼ 0:0002737
k1 ¼ 0:5000 g1 ¼ 1:134
and compute
h1 ¼ k1 k0 ¼ 0:5000 0:4670 ¼ 0:033

h2 ¼ k0 k2 ¼ 0:4670 0:4659 ¼ 0:0011
c ¼ h2 =h1 ¼ 0:0172=0:0292 ¼ 0:03333
and the coefficients

a ¼ ½cg1 g0 ð1 þ cÞ þ g2 = ch21 ð1 þ cÞ
¼ ½0:03333 ð1:134Þ 0:0002737 ð1 þ 0:03333Þ þ 0:04352

= 0:03333 0:0332 ð1 þ 0:03333Þ ¼ 145:1

b ¼ g1 g0 Ah21 =h1 ¼ ð1:134 0:0002737 145:1 0:0332 Þ=0:033 ¼ 39:16
c ¼ g0 ¼ 0:0002737
of the interpolating parabola. Then, k is estimated as follows

h
1=2 i
k ¼ k0 2c= b b2 4ac ¼ 0:4670
2 0:0002737=½39:16 ð39:162 4 145:1 0:0002737Þ1=2 ¼ 0:4670
Since this value of k is equal, within a tolerance of four significant figures, to the
value found in the preceding iteration, we accept it as the solution of g(k) = 0.
The optimum values of the mass ratios l1, l2, and l3 for the three stages result
from the following equations
l1 ¼ ðku1 1Þ=ðku1 e1 Þ ¼ ð0:4670 3:923 1Þ=ð0:4670 3:923 0:10Þ ¼ 4:542

l2 ¼ ðku2 1Þ=ðku2 e2 Þ ¼ ð0:4670 3:432 1Þ=ð0:4670 3:432 0:15Þ ¼ 2:507
l3 ¼ ðku3 1Þ=ðku3 e3 Þ ¼ ð0:4670 2:942 1Þ=ð0:4670 2:942 0:20Þ ¼ 1:361
The optimum values of the step masses m1, m2, and m3 are computed in
decreasing order by means of the following equations
m3 ¼ ðl3 1Þmpay =ð1 l3 e3 Þ ¼ ð1:361 1Þ 7000=ð1 1:361 0:20Þ ¼ 3472 kg

m2 ¼ ðl2 1Þ m3 þ mpay =ð1 l2 e2 Þ ¼ ð2:507 1Þ
ð3472 þ 7000Þ=ð1 2:507 0:15Þ ¼ 25290 kg

m1 ¼ ðl1 1Þ m3 þ m2 þ mpay =ð1 l1 e1 Þ ¼ ð4:542 1Þ
ð3472 þ 25293 þ 7000Þ=ð1 4:542 0:10Þ ¼ 232100 kg
The structural masses (ms)1, (ms)2, and (ms)3 of the stages result from
ðms Þ1 ¼ e1 m1 ¼ 0:10 232100 ¼ 23210 kg

ðms Þ2 ¼ e2 m2 ¼ 0:15 25290 3794 kg
ðms Þ3 ¼ e3 m3 ¼ 0:20 3472 694 kg
The masses of propellant (mprop)1, (mprop)2, and (mprop)3 contained in the stages
result from

mprop 1 ¼ m1 ðms Þ1 ¼ 232100 23210 208900 kg

mprop 2 ¼ m2 ðms Þ2 ¼ 25290 3794 21500 kg

mprop 3 ¼ m3 ðms Þ3 ¼ 3472 694 ¼ 2778 kg
The total mass of the rocket with the payload at lift-off is
m0 ¼ m1 þ m2 þ m3 þ mpay ¼ 232100 þ 25290 þ 3472 þ 7000 267900 kg
It remains to ascertain whether the value of m0 computed above for k = 0.4670

be indeed the minimum value. As has been shown above, this can be done by
evaluating the following expression
kui ð1 ei li Þ2 þ 2ei li 1
for i = 1, 2, 3. By so doing, there results
ku1 ð1 e1 l1 Þ2 þ 2e1 l1 1 ¼ 0:4670 3:923 ð1 0:10 4:542Þ2

þ 2 0:10 4:542 1 ¼ 0:4542
ku2 ð1 e2 l2 Þ2 þ 2e2 l2 1 ¼ 0:4670 3:342 ð1 0:15 2:507Þ2
þ 2 0:15 2:507 1 ¼ 0:3597
ku3 ð1 e3 l3 Þ2 þ 2e3 l3 1 ¼ 0:4670 2:942 ð1 0:20 1:361Þ2
þ 2 0:20 1:361 1 ¼ 0:2722
Since the values of preceding expression are greater than zero for i = 1, 2, 3, then
the value of m0 computed above for k = 0.4670 is indeed the minimum value.
11.4 Optimum Trajectory to Place a Satellite into Orbit
The present paragraph deals with the problem of determining the optimum trajec-
tory which a rocket must follow to place an artificial satellite into an orbit around
the Earth. The rocket is assumed to be a rigid body, and the aerodynamic forces
acting on it are neglected. In addition, the length of the powered part of the tra-
jectory is neglected in comparison with the radius of the Earth, and the central
gravitational field is considered as a field of parallel forces mg whose direction is
contained in an invariable plane xy, as shown in the following figure.
Let U be the angle which the thrust vector fT forms with the horizontal direction
at a time t. The gravity force vector mg, where m is the mass of the rocket, is
directed along the negative y-axis. The instantaneous velocity vector v is tangent to
the trajectory of the rocket. Let u and w be the components of the vector v along the
axes, respectively, x and y. The equations of motion are
du fT
¼ cos U
dt m
dw fT
¼ sin U g
dt m
dy
¼w
dt
11.4 Optimum Trajectory to Place a Satellite into Orbit 1283
The thrust per unit mass fT/m and the angle U vary with time. The trajectory of
the rocket depends on the laws of variation of these quantities with time. Supposing
to know how fT/m varies with time, we want to determine U as a function of time, in
order for the rocket to have the maximum horizontal velocity U at a given altitude
Y.
Let T be the time at which y reaches the value Y. This time depends on the
unknown function U = U(t) and therefore is not specified. However, it is reasonable
to assume that the propellant is exhausted before the time T.
The problem considered here differs from the problem of Sect. 11.3, where the
powered trajectory of the rocket has been assumed to be vertical. In the present
paragraph, we search the maximum value of an integral, whose integrand function
is specified by the first of the three equations written above, that is,
ZT
fT
U ¼ u0 þ cos U dt
m
0
where u0 is the horizontal component of the velocity vector v at the time t = 0 of

lift-off. The second and the third of these equations express constraints to which the
maximum value of the integral written above is subject. These constraints are
dw fT
sinU þ g ¼ 0
dt m
dy
w¼0
dt
As will be shown below, the horizontal velocity U, whose maximum value is

sought, can be rewritten as follows
ZT
fT dw fT dy
U ¼ u0 þ cos U þ k1 sin U þ g þ k2 w dt
m dt m dt
0
where k1 and k2 are two unknown functions of time.

With reference to the preceding figure, the conditions at the times t = 0 and t = T are
At t ¼ 0 At t ¼ T
x¼0
y ¼ y0 y¼Y
u0 ¼ v0 cos h0
w0 ¼ v0 sin h0 w¼0
The magnitude v0 of initial velocity vector v0 is fixed, but the initial angle h0 is
not determined, and therefore, the term u0 in front of the integral written above
gives a contribution to the horizontal velocity U. The term u0 = v0 cos h0 gives a

contribution −v0 sin h0 dh0 to the variation dU of the horizontal velocity U.
For a better understanding of the equations written above, it is necessary to
discuss some fundamental concepts of variational calculus. In general terms, the
problem considered in the present paragraph consists in finding the maximum or the
minimum value of an integral I of the following type
Zb
dz
I ¼ f t; z; dt
dt
a
where the integral I is taken along a curve z = U(t), f is a functional, that is, a function
having functions as its arguments and returning a scalar, and z and dz/dt designate,
respectively, one or more functions and its or their first derivatives with respect to the
independent variable t.
Since the value of the integral I depends on the curve z = U(t), then the problem
to be solved consists in determining the unknown optimum function z = U(t) such
that the integral I takes its maximum or minimum value.
Let z = U(t) be the optimum curve going from a to b. Let z1 = U1(t) be another
(non-optimum) curve going from a to b1 (where b1 6¼ b). The quantity z computed
along the curve z1 = U1(t) can be expressed as follows
z1 ¼ z þ dz
where dz is the variation of z. The variation dz differs from the differential dz,
because dz is the increment computed along the curve z due to an increment dt of
time t, whereas dz is the difference between the two curves z and z1 at any given
time t. The difference dz′ between the slopes of the curves z1 and z at a time t is
defined as follows
dz0 ¼ z01 z0
Assuming dz to be a continuous function of time, the expression dz = z1 − z can

be differentiated as follows
dðdzÞ
¼ z01 z0
dt
By comparing the two equations written above, there results

dðdzÞ dz
¼d
dt dt
In other words, the differentiation (d/dt) and the variation (d) are two operations
R
which can be executed in any order. The same rule applies to the integration ( ) and
the variation (d), that is,
Z Z
dzdt ¼ d zdt
The variation of the functional under the integral sign can be expressed along the
curve z1 by means of an expansion in Taylor series about the curve z, as follows
@f @f
f ðt; z þ dz; z0 þ dz0 Þ ¼ f ðt; z; z0 Þ þ dz þ 0 dz0 þ
@z @z
Considering only the first-order variation, the variation of the integral I is
Zb
@f @f 0
dI ¼ dz þ 0 dz dt
@z @z
a
The second term in parentheses can be integrated by parts. This yields
b Zb
@f @f d @f
dI ¼ dz þ dz dt
@z0 a @z dt @z0
a
In order for the curve z = U(t) to be the optimum curve which makes the value
of I minimum or maximum, the variation dI written above must be equal to zero.
In many cases, the endpoints a and b of the curve z = U(t) coincide with the
endpoints a and b1 of the curve z1 = U1(t), which means b1 = b. In such cases, the
first term on the right-hand side of the preceding equation is equal to zero.
Since the variation dz under the sign of integral is arbitrary and is therefore not
necessarily equal to zero, then the expression in parentheses must be equal to zero
in order for dI to be equal to zero. This implies
@f d @f
¼0
@z dt @z0
which is the well known Euler–Lagrange equation. When this equation is satisfied,
then the value of I is either maximum or minimum.
When the variation dz is not equal to zero at the endpoints, then the condition to
be satisfied is
b Zb
@f @f d @f
dI ¼ dz þ dzdt ¼ 0
@z0 a @z dt @z0
a
What has been shown above holds in case of a search for the minimum or the
maximum value of I in the absence of constraints.
By contrast, in the presence of constraints expressed by the following equations

dz
g1 t; z; ¼0
dt

dz
g2 t; z; ¼0
dt
..
.
it is necessary to find the minimum or the maximum value of the following integral
Zb
dz dz dz
I¼ f t; z; þ k1 g1 t; z; þ k2 g2 t; z; þ dt
dt dt dt
a
This assures the satisfaction of these constraints. The quantities k1, k2, … are
also functions of the variables which appear in the problem.
The preceding equation
b Zb
@f @f d @f
dI ¼ dz þ dz dt ¼ 0
@z0 a @z dt @z0
a
can also be applied in the presence of constraints, but in this case the function
f appearing above becomes

dz dz dz
f ¼ f t; z; þ k1 g1 t; z; þ k2 g2 t; z; þ
dt dt dt
Now, coming back to the particular problem described in this paragraph, the
variation dU of U can be expressed as follows
ZT
fT dw dy
dU ¼ v0 sin h0 dh0 þ d f ; U; w; ; dt
m dt dt
0
where the variation dI of the integral written above is
b Zb
@f @f d @f
@z0 a @z dt @z0
a
The acceleration fT/m due to the thrust acting on the rocket is assumed to be a
known function of time. The unknown function of time to be determined is U(t),
that is, the optimum function which makes the value of U maximum at a given
altitude Y. Now, we evaluate the partial derivatives ∂f/∂z and ∂f/∂z′ which appear in
the preceding equation. Taking account of
ZT
fT dw fT dy
m dt m dt
0
these partial derivatives are

@f fT @f
¼ ðsinU þ k1 cosUÞ ¼0
@/ m @/0
@f @f
¼ k2 ¼ k1
@w @w0
@f @f
¼0 ¼ k2
@y @y0
These partial derivatives are inserted into

b Zb
@f @f d @f
@z0 a @z dt @z0
a
By so doing, the total variation dU results
dU ¼ v0 sin h0 dh0 þ ½k1 dwT0 þ ½k2 dyT0

ZT
fT dk2 dk1
ðsin U þ k1 cos UÞdU þ dy þ þ k2 dw dt
m dt dt
0
Since y has fixed values (y0 and Y) at the endpoints, respectively, 0 and T, then
the variation dy at these endpoints is equal to zero. Therefore, ½k2 dyT0 = 0.
As to the term ½k2 dwT0 , the vertical component of the velocity vector v at t = 0 is
w0 ¼ v0 sin h0
where v0 has a fixed value. Therefore
½dw0 ¼ d½v0 sin h0 ¼ v0 cos h0 dh0
At t = T, the vertical component of the velocity vector v results from integrating

the following expression from 0 to T
dw fT
sin U þ g ¼ 0
dt m
This integration yields
ZT ZT
fT fT
w¼ sin U g dt ¼ sin U dt gT þ v0 sin h0
m m
0 0
The variation of w at t = T, due to the variations dU and dh0, is
ZT
T fT
½dw ¼ cos UdU dt þ v0 cos h0 dh0
m
0
For the optimum curve, w becomes equal to zero at t = T. For another curve,
w becomes equal to zero at another time t = T + dT.
Since fT/m is not varied and is equal to zero for t > T, then
ZT
T fT
½dw ¼ cos UdU dt þ v0 cosh0 dh0 ¼ gdT
m
0
In other words
½dwT0 ¼ gdT v0 cos h0 dh0
and then
½k1 dwT0 ¼ gk1;t¼T dT v0 k1;t¼0 cos h0 dh0
Therefore, the preceding equation expressing the variation of U
dU ¼ v0 sin h0 dU0 þ ½k1 dw T0 þ ½k2 dy T0

ZT
fT dk2 dk1
m dt dt
0

dU ¼ v0 sin h0 þ k1;t¼0 cos h0 dh0 þ gk1;t¼T dT
ZT
fT dk2 dk1
m dt dt
0
Since the variations dh0, dT, dU, dy, and dw are arbitrary, then the preceding
expression of dU can be equal to zero only if the coefficients which multiply these
variations are, all of them, equal to zero. By setting the coefficient of dh0 equal to
zero and simplifying, there results
tan h0 ¼ k1;t¼0
By doing the same with the coefficients of dT, dU, dy, and dw, there results
dk2 dk1
k1;t¼T ¼ 0 tan U ¼ k1 ¼0 þ k2 ¼ 0
dt dt
The equation dk2/dt = 0 written above implies
k2 ¼ constant ¼ c2
This result, substituted into the equation dk1/dt + k2 = 0, yields
dk1
c2 ¼ 0
dt
which, integrated, yields
k1 ¼ c 2 t þ c 1
At t = 0, there results k1, t=0 = c1 = −tan h0 because tan h0 = −k1, t=0.

At t = T, there results k1, t=T = 0 = c2T + c1 = c2T − tan h0. Therefore,
c1 = −tan h0 and c2 = (tan h0)/T. These results, substituted into k1 = c2t + c1, yield
t
k1 ¼ 1 tan h0
T
Since tan U = −k1, then the required optimum function U = U(t) is

h t i
U ¼ arctan 1 tan h0
T
The optimum function U = U(t) found above takes the value h0 at t = 0 and the
value 0 at t = T. The method described in the present paragraph is due to
Okhotsimskii and Eneev [10].
11.5 Optimum Consumption of Propellant
In Sect. 11.4, we found the optimum function U = U(t) which corresponds to the
maximum value of the horizontal component U of the velocity of a rocket at the end
of the propelled part of its trajectory, without reference to the manner in which the
thrust per unit mass fT/m varies with time. In that case, the initial angle h0, which
the velocity vector v0 of the rocket forms with the horizontal direction, was allowed
to vary. The optimum function U = U(t) was found to be
t
tan U ¼ 1 tan h0
T
If we consider the initial angle h0 as a fixed quantity, that is, as a quantity which
is not allowed to vary, then tan U is a linear function of time, which can be
tan U ¼ tan U0 c2 t
In either case, each of the two functions U = U(t) defined above corresponds to
the maximum value of the horizontal component of the velocity vector v at t = T for
the specified function fT/m. However, the maximum value of the horizontal com-
ponent U of the velocity vector v at t = T depends on the chosen function fT/m and
consequently on the consumption of propellant as a function of time. In the present
paragraph, we want to determine the optimum function of time expressing the
propellant consumption, which corresponds to the largest of these values of U, all of
them corresponding to the optimum function U = U(t).
The preceding figure shows a possible graph of the velocity V of a rocket as a

function of time, where Vk is the maximum velocity which can theoretically be
reached in the absence of all forces except thrust. Vk is a function of the specific
impulse Isp of the propellant and also of the mass ratio l = m0/mb0 of the rocket. As
has been shown in Sect. 11.1, this function is

m0
Vk ¼ g0 Isp ln
mb0
where the value of Isp is assumed to be independent of the consumption of

propellant.
11.5 Optimum Consumption of Propellant 1291
If we set
dV fT
¼
dt m
that is, if the thrust per unit mass is the sole cause of acceleration of a rocket, then
the velocity V V(t) of the rocket can be expressed as a function of time in the
following manner
Zt
fT
V ðtÞ ¼ ds
m
0
where V is limited by the maximum theoretical value Vk. The velocity V depends on
time and also on the manner in which the propellant is consumed. Any curve having
positive slope V′ dV/dt is a possible curve for the function V(t). Since fT and
m have positive values, the slope dV/dt may vary between zero (if fT = 0) and
infinity (if the propellant were consumed instantaneously).
Remembering the following equation of Sect. 11.4
ZT
fT dw fT dy
m dt m dt
0
and expressing the thrust per unit mass as follows
fT dV
¼
m dt
the integral, whose maximum value is sought in the present case, is
ZT
dV dw dy
U¼ ðcos U k1 sin UÞ þ k1 þ g þ k2 w dt
dt dt dt
0
Since the optimum function U = U(t) expressing the thrust angle is used in the
present case, then all of the possible variations have been considered in the problem
of Sect. 11.4 except dV. Therefore, it is necessary to consider only the first term of
the integral written above, that is,
ZT
dV
U¼ ðcos U k1 sin UÞ dt
dt
0
In order to determine the terms which appear in the general equation
b Zb
@f @f d @f
dI ¼ dz þ dzdt ¼ 0
@z0 a @z dt @z0
a
we evaluate the partial derivatives ∂f/∂z and ∂f/∂z′ as follows
@f @f
¼0 ¼ cos U k1 sin U
@V @V 0
Since dV is equal to zero at the endpoints t = 0 and t = T, then the variation dU

of
ZT
dV
U¼ ðcos U k1 sin UÞ dt
dt
0
is
ZT
d
dU ¼ ðcos U k1 sin UÞdVdt
dt
0
ZT
dU dk1
¼ ðsin U þ k1 cos UÞ þ sin U dVdt
dt dt
0
Since
sin U þ k1 cos U ¼ 0
then
k1 ¼ tan U
Since tan U is a linear function of time, that is, tan U = tan U0 − c2t, then
k1 ¼ tan U ¼ tan U0 þ c2 t
Therefore
dk1 d
¼ ð tan U0 þ c2 tÞ ¼ c2
dt dt
and the equation
ZT
dU dk1
dU ¼ ðsin U þ k1 cos UÞ þ sin U dVdt
dt dt
0
ZT
dU ¼ c2 ðsin UÞdVdt
0
Remembering the trigonometric identity
tan a
sin a ¼ 1
ð1 þ tan2 aÞ2
and the equation
tan U ¼ tan U0 c2 t
there results
8 9
ZT ZT >
< >
=
tan U0 c2 t
dU ¼ c2 ðsin UÞdVdt ¼ c2 h i1 dV dt
: 1 þ ðtan U c tÞ2 2 >
> ;
0 0 0 2
The preceding expression of dU must be equal to zero, in order to determine the

optimum function of time expressing the propellant consumption. Let us set for
convenience
tan U0 c2 t
W¼h i12
1 þ ðtan U0 c2 tÞ2
so that the condition dU = 0 becomes
ZT
dU ¼ c2 WdVdt ¼ 0
0
The function W specified above depends on time t and may take positive or
negative values as t varies from 0 to T. Consequently, in order for the integral
written above to be equal to zero, the function V = V(t) must be a discontinuous

function of time, that is, the fuel must be burned instantaneously.
Let us consider the following three cases.
(1) The function W is always positive within the interval 0 t T. In this case,
the fuel is burned instantaneously at t = 0, as shown in the following figure.
The optimum function V = V(t) is the discontinuous function shown above. Since
V = V(t) is the optimum, dV and dU would be negative for all variation of V from
its optimum value. The variation dV of V from the optimum value cannot be
positive, because V cannot be greater than Vk, which in turn depends on the
specific impulse Isp and also on the mass ratio m0/mb0 of the rocket, according to

m0
Vk ¼ g0 Isp ln
mb0
(2) The function W is always negative within the interval 0 t T. In this case,
the fuel is burned instantaneously at t = T, as shown in the following figure.
The variations dV and dU from their optimum values would be both of them
positive.
(3) The function W changes sign from plus to minus within the interval 0 t T. In
this case, part of the fuel is burned instantaneously at t = 0, and the remaining part
is burned instantaneously at t = T, as shown in the following figure.
Any varied curve V(t) whose slope is zero or positive is allowable. If a positive
dU cannot be produced by using this curve, the discontinuous curve shown above is
the optimum curve. The method described in the present paragraph is due to
Thomson [2].
11.6 Gravity Turn Trajectories
A gravity turn (or zero-angle-of-attack, or zero-lift) trajectory is a particular path

which uses the force of gravity to insert a spacecraft into an orbit around the Earth.
The same path can also be used to de-orbit a spacecraft which is meant to land to a
celestial body, which may be a planet of destination or the Moon.
This particular trajectory is obtained by simply keeping the propulsive thrust fT
of a spacecraft always aligned with the velocity vector v of the spacecraft itself [11].
No attempt is made to reduce the consumption of propellant to a minimum; instead,
we seek here the minimum value of the stresses acting in a direction perpendicular
to the longitudinal axis of the spacecraft, because the thrust fT and the aerodynamic
drag fD act only along this axis.
A space vehicle which moves along a gravity turn trajectory does not use the
thrust applied to it to change its direction of flight, and therefore, the maximum
possible part of the thrust is used to accelerate the vehicle into the desired orbit. In
addition, the same vehicle can maintain a low or even zero value of angle of attack,
which fact reduces the transverse aerodynamic stresses due to lift.
For this purpose, rockets which are aerodynamically stabilised use fixed fins.
Instead, mechanically stabilised rockets use control systems to generate moments
apt to reduce their angles of attack [12], as will be shown below.
As has been done in the preceding paragraphs, so here the length of the powered
part of the trajectory is neglected in comparison with the radius of the Earth, and the
central gravitational field is considered as a field of parallel forces whose direction
is contained in the invariable plane xy shown in the preceding figure, where x is the
downrange distance measured on the surface of a flat Earth, and y is the altitude of
the centre of mass O of the rocket over this surface.
Let w be the angle which the velocity vector v of the rocket forms with the
y-axis, that is, with the local vertical, at a time t. Neglecting the aerodynamic drag
force, the two scalar equations of motion of the rocket along the axes x and y are
dvx vx
m ¼ fT sin w ¼ fT
dt v
dvy vy
m ¼ fT cos w mg ¼ fT mg
dt v
11.6 Gravity Turn Trajectories 1297
where
12
v ¼ v2x þ v2y
is the magnitude of the velocity vector v of the rocket.

By multiplying all terms of the first equation by vx, all terms of the second
equation by vy, and summing, there results

dvx dvy
m vx þ vy ¼ fT v mgvy
dt dt
In addition, by differentiating v2 = v2x + v2y with respect to time, there results
dv dvx dvy
2v ¼ 2vx þ 2vy
dt dt dt
Therefore, the first of the two equations written above becomes
dv fT vy f T
¼ g ¼ g cos w
dt m v m
Likewise, by multiplying all terms of the first of the same two equations written
above by vy, all terms of the second equation by vx, and subtracting, there results

dvx dvy
m vx vy ¼ mgvx
dt dt
Since

dðtan wÞ d vx 1 dvx dvy
¼ ¼ 2 vx vy
dt dt vy vy dt dt
then
dðtan wÞ 1 dw
v2y ¼ v2y ¼ gvx
dt cos2 w dt
and therefore, the second of the two equations written above becomes
dw vx 1 g
¼g cos2 w ¼ sin w
dt vy vy v
This is because vx/vy = tan w, and vy = v cos w. Summarising, the two scalar
equations of motion for a rocket which moves along a gravity turn trajectory are
dv fT
¼ g cos w
dt m
dw g
¼ sin w
dt v
Some authors (see, e.g., Ref. [12]) write these equations in terms of the velocity
v and the flight path angle c = p/2 − w. In this case, the two equations written above
become
dv fT
¼ g sin c
dt m
dc g
¼ cos c
dt v
These differential equations are nonlinear. No analytical solution is known in the

general case, that is, when fT/(mg) varies arbitrarily with time. They can be inte-
grated numerically (by using, e.g., one of the methods shown in Chap. 6), when the
thrust fT and the mass m of the rocket are known functions of time. In order to
integrate analytically these equations, Thomson [2] assumes fT/(mg) to be constant
in time. This happens when:
• the mass m decreases linearly with time, because of the consumption of pro-
pellant, as follows
m ¼ m0 m0 ðt t0 Þ
where the mass flow m′ is constant, and

• the thrust fT also decreases linearly with time, in order to conform with the
decreasing mass m of the rocket.
The acceleration of gravity g is assumed to be constant
(g g0 = 9.80665 m/s2).
By so doing, the variable quantity fT/(mg) is assumed to be piecewise constant in
time, that is, constant within small intervals of time. Therefore, a numerical inte-
gration can be carried out by using the analytical solution which can be found
within each of these intervals.
In the sequel, the notation is simplified by setting fT/(mg) = n = constant within
each interval. By defining a new variable

w
z ¼ tan
2
and remembering the following identity

1
w 1 cos w 2 sin w
z ¼ tan ¼ ¼
2 1 þ cos w 1 þ cos w
there results
1 cos w
z2 ¼
1 þ cos w
which, solved for cos w, yields
1 z2
cos w ¼
1 þ z2
By differentiating z = tan (w/2) with respect to time, there results

dz 1 d w 1 dw
¼ ¼
dt w dt 2 w dt
cos2 2cos2
2 2
and therefore

dw w dz 1 þ cos w dz dz
¼ 2 cos2 ¼2 ¼ ð1 þ cos wÞ
dt 2 dt 2 dt dt
The expressions cos w = (1 − z2)/(1 + z2) and dw/dt = (1 + cos w)(dz/dt)

derived above and also n = fT/(mg) are substituted into the two differential
equations
dv fT
¼ g cos w
dt m
dw g
¼ sin w
dt v
This substitution yields
1 dv 1 z2
¼n
g dt 1 þ z2
v dz
¼z
g dt
The two equations written above may be combined together by eliminating

dt between them. This yields
dv dz 1 z2 dz
¼n
v z 1 þ z2 z
By integrating the preceding differential equation, there results

Z Z Z
1 z2 dz z dz
lnðvÞ ¼ lnðzn Þ dz ¼ lnðzn Þ þ
zð1 þ z2 Þ zð1 þ z2 Þ 1 þ z2
The two integrals on the right-hand side of the preceding equation can be
computed (omitting the constant of integration) as follows
Z Z Z
zdz 1 2zdz 1 dð 1 þ z 2 Þ 1

¼ ¼ ¼ ln 1 þ z2
1þz 2 2 1þz 2 2 1þz 2 2
Z Z 2
dz 1 z 1
1 z
¼ dz ¼ ln ð zÞ ln 1 þ z 2
¼ ln
zð1 þ z Þ
2 z 1þz 2 2 2 1 þ z2
Therefore, taking account of the constant of integration expressed in the form of

ln(c), the solution of the preceding differential equation is
2
1 z 1

lnðvÞ ¼ lnðz Þ ln
n
þ ln 1 þ z2 þ lnðcÞ
2 1þz 2 2

1 þ z 2
¼ lnðzn Þ þ ln þ lnðcÞ
z

v ¼ czn1 1 þ z2
In order to evaluate c, we impose the initial conditions, as follows

v0 ¼ czn1
0 1 þ z20
hence
v
c¼
0 2
zn1
0 1 þ z0
where z0 = tan (w0/2). By substituting

v ¼ czn1 1 þ z2
into
v dz
¼z
g dt
and then integrating from z0 to z, there results
Zz Zz Zz z
c
c c c 1 z2
t¼ zn2 1 þ z2 dz ¼ zn2 dz þ zn dz ¼ zn1 þ
g g g g n 1 n þ 1 z0
z0 z0 z0

c 1 z2 1 z20
¼ zn1 þ zn1 þ
g n 1 nþ1 0
n 1 nþ1
The preceding equation expresses t as a function of z = tan (w/2) for any con-
stant c. Therefore, w = 2 arctan(z) is known as a function of t for any c.
The velocity v depends on z and c through v = czn−1(1 + z2), and therefore, v is
also known as a function of t and c. Summarising, the following equations

v ¼ czn1 1 þ z2
v0
c ¼ n1

z0 1 þ z20

c n1 1 z2 1 z20
t¼ z þ z0
n1
þ
g n 1 nþ1 n 1 nþ1
yield the solution for the gravity turn trajectory when fT/(mg) = n = constant.
In order to apply these equation to the case of fT/(mg) varying in time, we start
from the initial conditions expressed by z0 = tan (w0/2), c = v0/[zn−1
0 (1 + z20)], and
fT/(mg) = n at t = t0.
Then, we choose a value of w slightly greater than w0 (or an increment Dw of w0)
and determine v by means of v = czn−1 (1 + z2) and Dt by means of

c n1 1 z2 1 z20
Dt ¼ z þ z0
n1
þ
g n 1 nþ1 n 1 nþ1
The increment in the displacement of the rocket along the gravity turn trajectory
is computed as follows
1
Dx ¼ ðv0 sin w0 þ v sin wÞDt
2
1
Dy ¼ ðv0 cos w0 þ vcos wÞDt
2
Then, this procedure is repeated with the values relating to the new point taken
as initial conditions.
The procedure described above is due to Thomson [2]. More recently, Sharaf
and Alaqal [13] have proposed a numerical algorithm based on this procedure. This
algorithm computes the co-ordinates x and y and the tangential velocity v of a space
vehicle which moves along a gravity turn trajectory with varying thrust-to-weight
ratio. The particulars on the algorithm can be found in Ref. [13].
In case of multi-stage rockets, the thrust fT and the mass m of the rocket vary
discontinuously in time. A discontinuity in thrust and mass occurs at the moment in
which each extinguished stage separates from the remaining part of the rocket.
Therefore, in case of a multi-stage rocket moving along a gravity turn trajectory, the
whole procedure described above is to be executed more than once, depending on
the number of stages fired during this trajectory.
In practice, as has been shown in Sect. 4.5, a rocket is launched vertically and
then pitched over after a few seconds, in order to steer its longitudinal axis towards
the downrange direction and to put it in the desired orbital plane. The pitch-over
manoeuvre is performed by gimballing slightly the rocket engines to one side, in
order for a torque to act on the rocket. Therefore, a small component of the force of
gravity acts perpendicularly to the longitudinal axis of the rocket. From this
moment until the moment of insertion into orbit, this perpendicular component
increases and causes the velocity vector v of the rocket to rotate towards the
horizontal direction during ascent. The initial pitch-over angle w0 depends on the
specific rocket used and also on the desired orbit. After completion of the pitch-over
manoeuvre, the rocket engines are oriented again in their previous non-gimballed
direction.
11.7 Trajectories of Long-Range Ballistic Missiles
By ballistic missile, we mean any missile which does not rely upon aerodynamic
surfaces to produce lift. In the absence of these surfaces, lift is negligible in
comparison with the other forces acting upon the missile, which are essentially
thrust (fT), drag (fD), and gravity (fG = mg). After the thrust force is terminated, a
ballistic missile is acted upon only by gravity and aerodynamic drag.
In order to take account of drag, the following equations of motion derived in
Sect. 11.6
dv fT
¼ g cos w
dt m
dw g
¼ sin w
dt v
11.7 Trajectories of Long-Range Ballistic Missiles 1303
are modified here as follows
dv fT fD
¼ g cos w
dt m
dw g
¼ sin w
dt v
where v is the magnitude of the velocity vector v of the missile, m is the instan-
taneous mass of the missile, g g0 = 9.80665 m/s2 is the magnitude of the
acceleration of gravity, and w is the angle which the velocity vector v of the missile
forms, at a time t, with the local vertical.
dv dv dw dv g f f
T D
¼ ¼ sin w ¼ g cos w
dt dw dt dw v m
which yields after simplification
1 dv fT fD
¼ cot w
v dw mg sin w
The preceding differential equation governs the motion of a ballistic missile, in

accordance with Ref. [2]. Of course, the term fT takes the value zero after the time tb
of burnout.
In particular, long-range ballistic missiles are ballistic missiles whose range

(usually greater than 3500 km) is comparable with the mean radius of the Earth
(rE = 6371 km). In this case, the shortest distance between two points on the surface
of a spherical Earth is to be measured along a great circle, and the trajectory
followed by a long-range ballistic missile lies on the same plane as that of the great
circle, as shown in the preceding figure.
As has been shown in Sect. 1.1, the trajectory of the missile is an ellipse having
its occupied focus in the centre of mass of the Earth, its perigee inside the Earth, and
its apogee (placed at a distance ra from the occupied focus) coinciding with the
point of maximum height above the surface of the Earth.
We want to determine the range (rEh), the maximum height (h), and the time of
flight (Dt) of the missile, as will be shown below.
Let r0, v0, and c0 be, respectively, the radius vector, the velocity, and the flight
path angle of the missile at the time tb of burnout. As has been shown in Sect. 4.6,
the true anomaly /0 of the missile at the time tb and the eccentricity e0 of its elliptic
trajectory result from

2
tan /0 ¼

r0 v20 =lE cos2 c0 1
2 2
r 0 v0
e ¼
2
1 cos2 c0 þ sin2 c0
lE
where lE is the gravitational parameter of the Earth. Let h be the angle between the
radius vector at the time tb of burnout and the radius vector at the final time tf in
which the missile hits the surface of the Earth.
By inspection of the preceding figure, there results
h
¼ p /0
2
and consequently

h h
tan /0 ¼ tan p ¼ tan
2 2
Therefore, the equation

2
tan /0 ¼
2
r0 v0 =lE cos2 c0 1

2
h r0 v0 =lE sin c0 cos c0
tan ¼
2
2 r0 v0 =lE cos2 c0 1
The preceding equation makes it possible to determine h (and therefore the range
rEh) as a function of c0 for any given couple of values of r0 and v0.
In order to plot h versus c0, we set for convenience
r0 v20
k
lE
and consider k as a parameter, to which we give arbitrarily the following values:

0.50, 0.75, 1.00, 1.25, 1.50, 1.75, and 2.00. The resulting plot is shown in the
following figure.
After h has been determined as a function of r0, v0, and c0 as has been shown
above, the maximum height h reached by the missile over the surface of the Earth
results from
h ¼ r a r E ¼ að 1 þ e Þ r E
where ra is the radius vector corresponding to the apogee of the elliptic trajectory of
the missile, a is the major semi-axis of this trajectory, and rE is the mean radius of
the Earth.
By dividing all terms of the preceding equation h = a(1 + e) − rE by rE, there

results
h a
¼ ð 1 þ eÞ 1
rE rE
On the other hand, we observe that, for /0 = p – h/2, there results r0 = rE.
Remembering the equation of an elliptic orbit
að 1 e 2 Þ
r¼
1 þ e cos /
and replacing r with rE and / with p − h/2, there results

a 1 e2 a 1 e2
rE ¼ ¼
h h
1 þ e cos p 1 e cos
2 2
and therefore

h
1 e cos
a 2
¼
rE 1 e2
By substituting the preceding expression of a/rE into
h a
¼ ð 1 þ eÞ 1
rE rE
there results after simplification

h e h
¼ 1 cos
rE 1 e 2
which determines the maximum height h reached by the missile.

It remains to determine the time of flight Dt = tf − tb, where tf is the time in
which the missile hits the surface of the Earth, and tb is the time of burnout.
For this purpose, the equation h = a(1 + e) − rE solved for a yields the value of
the major semi-axis of the elliptic trajectory, as follows
h þ rE
a¼
1þe
With reference to the following figure, let us consider the triangle whose vertices
are the centre of mass of the Earth (O), the point of burnout (P0) of the missile, and
the apogee of the elliptic trajectory. The sides of this triangle are ra = rE + h, rE, and c.
The angle formed by the sides rE and ra is h/2.
Since an elliptic orbit is symmetrical with respect of its line of apsides, then the time
of flight Dt = tf − tb from P0 to the point in which the missile hits the surface of the
Earth is equal to two times the time of flight from P0 to the apogee of the elliptic
trajectory. The true anomalies of the point P0 and of the point of apogee are,
respectively, /0 = p − h/2 and /a = p. We compute Dt = tf − tb by means of a
formula, due to Cunningham [14], which has been shown in Sect. 3.17 and is
rewritten below for convenience of the reader.
12 h i
a3
1
Dt ¼ 2 e 1 e2 2 ðp0 pa Þ þ 2ðqa q0 Þ
lE
where

h
sin
sin /0 2
p0 ¼ ¼
1 þ e cos /0 h
1 e cos
2
sin /a sin p
pa ¼ ¼ ¼0
1 þ e cos /a 1 þ e cos p
" 1 # " 1 #
1e 2 /0 1e 2 h
q0 ¼ arctan tan ¼ arctan cot
1þe 2 1þe 4
" 1 # " 1 #
1e 2 /a 1 e 2 p p
qa ¼ arctan tan ¼ arctan tan ¼
1þe 2 1þe 2 2
The preceding equations can be combined together into one as follows

12 ( " 1 #)
a3
12 sin /0 1e 2 /
Dt ¼ 2 e 1e 2
þ p 2 arctan tan 0
lE 1 þ e cos /0 1þe 2
which is just the equation indicated by Thomson [2].
References
1. R.W. Buchheim et al., Space Handbook: Astronautics and Its Applications (United States
conghand/propelnt.htm
2. W.T. Thomson, Introduction to Space Dynamics (Dover Publications, New York, U.S.A,
1986). ISBN 0-486-65113-4
3. Smithsonian National Air and Space Museum, Saturn V: America’s Moon Rocket, Space
Race, Racing to the Moon. Web site https://airandspace.si.edu/exhibitions/space-race/online/
sec300/sec384.htm
4. NASA, Space Flight Systems, Glenn Research Centre, Booster Staging. Web site https://
spaceflightsystems.grc.nasa.gov/education/rocket/rktstage.html
5. NASA, Multi-stage and clustered rockets, 15 p., Glenn Research Centre Explorers Posts. Web
site https://explorersposts.grc.nasa.gov/post630/08-09%20Files/Rocket%20Design%
20Mission%20Discussion/Multi-Stage
6. D.H. Mitchell et al., Flight separation mechanisms, NASA SP-8056, October 1970, 38 p.,
monograph available at the web site http://ntrs.nasa.gov/archive/nasa/casi.ntrs.nasa.gov/
19710019510.pdf
2005). ISBN 0-7506-6169-0
8. J.J. Coleman, Optimum stage-weight distribution of multistage rockets, Space Technology
Laboratories, Inc., STL/TN-60-0000-09036, 23 pages, 3 February 1960. Web site http://www.
dtic.mil/dtic/tr/fulltext/u2/607412.pdf
References 1309
9. H.H. Hall, E.D. Zambelli, On the optimization of multistage rockets. J. Jet Propuls. 28(7),
463–465 (1958)
10. D.E. Okhotsimskii, T.M. Eneev, Some variation problems connected with the launching of
artificial satellites of the Earth. J. Br. Interplanet. Soc. 16(5), 263–294 (1958)
11. G.J. Culler, B.D. Fried, Universal gravity turn trajectories. J. Appl. Phys. 28(6), 672–676
(1957)
12. J.W. Cornelisse, H.F.R. Schöyer, K.F. Wakker, Rocket Propulsion and Spaceflight Dynamics
(Pitman, London, 1979). ISBN 0-273-01141-3
13. M.A. Sharaf, L.A. Alaqal, Computational algorithm for gravity turn maneuver. Glob. J. Sci.
Front. Res. Math. Decis. Sci. 12(13), Version 1 (2012), Global Journals Inc. (USA), 7 p., web
site https://globaljournals.org/GJSFR_Volume12/6-Computational-Algorithm-for-Gravity.pdf
14. F.G. Cunningham, Calculation of the eclipse factor for elliptical satellite orbits, NASA TN
D-1347, June 1962

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Springer Aerospace Technology

Alessandro de Iaco Veris

ISSN 1869-1730 ISSN 1869-1749 (electronic)

© Springer International Publishing AG 2018

Printed on acid-free paper

This Springer imprint is published by Springer Nature

Ingegnati, se puoi, d'esser palese

(Dante, Vita nova, XIX).

Rome, Italy Alessandro de Iaco Veris

1 The Two-Body Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1

2.8 Improvement in Orbit Determination by Differential

3.22 Atmospheric Density Models . . . . . . . . . . . . . . . . . . ........ 522

5.12 The Aerodynamic Assist . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 838

7.11 Elementary Concepts on Elliptic Integrals . . . . . . . . . . . . . . . . 1025

9.15 Three-Axis Controlled Satellites . . . . . . . . . . . . . . . . . . . . . . . . 1203

1.1 Position of the Problem

Astrodynamics is deﬁned by Kaplan [1] as “the study of controlled flight paths of

© Springer International Publishing AG 2018 1

where l = GM = 1.327 1011 km3/s2 is the gravitational parameter of the Sun,

where a = dv/dt is the acceleration vector applied to the body.

and shown in the following ﬁgure.

This is because, by deﬁnition, r = r1 − r2. Likewise, there results

The attracting force f1 is a Newtonian gravitational force, such that

Because of the identity r r 0, the preceding expression reduces to

Since (r r′)′ = r′ r′ + r r″ = r r″, and since r′ r′ = 0, r r″ = 0,

Thus, the instantaneous areal velocity of the planet is

By applying the following identity (rule of triple vector products)

to the expression −(l/r3)r (r r′), there results

Since r  r′ = rr′ and r  r = r2, then

Since h is a constant vector, then the preceding differential equation can be

where e is a constant of integration called the eccentricity vector. The preceding

is equal to 0. This is because (r′ h) is orthogonal to h. Therefore, the direction of

As to the right-hand side of the same equality, there results

Thus, the equality can also be written as follows

which, solved for r, yields

The expression written above is the equation of a conic section in polar

In this connexion, it is to be remembered that the semi-latus rectum of a conic

has important consequences for the mechanical energy possessed by a planet

As has been shown above, there results

This equation, solved for e, yields

Since (r′ h)  (r′ h) = r′  h (r′ h) and since r′ and h are orthogonal

On the other hand,

The preceding equation, remembering that p = h2/l, becomes

where the quantity a has the dimension of length, there results

The previous expression, inserted into 1 − e2 = p(2/r − v2/l), leads to

which is the so-called energy integral or vis-viva integral.

Since p is never negative, then e is less or greater than unity, depending on

1.2 The Conic Sections and Their Geometrical Properties

is a conic section having the following equation in polar co-ordinates (r, /)

According to the most common deﬁnition, a conic section is a curve resulting

The generation of a conic surface may be described as follows. Let us consider a

Let e, called eccentricity, be the positive constant of proportionality. If e = 1,

where p is the semi-latus rectum, which is shown in the preceding ﬁgure.

In case of a parabola (e = 1), the general equation of conic sections in polar

Since the following expression

is a trigonometric identity, then

Since Sect. 1.1 has shown that

into the preceding equation yields

By simplifying and separating the variables, the preceding differential equation

is a trigonometric identity, then the preceding differential equation becomes

which, integrated, yields

then, integrating by parts, there results

is a trigonometric identity, then

By deﬁning the quantity B = 3(l/p3)½ (t − tp), the preceding equation reduces to

Since r r′ = rr′ and r r = r2, then

Since (r′ h) (r′ h) = r′ h (r′ h) and since r′ and h are orthogonal

x2 y2 1 sin2 f 1 sin2 f