MATHS TOPIC

1. Basic coordinate
2. Straight Lines
3. Simultaneous Equation
4. Polynomial
5. Algebra of polynomial
6. Matrices
7. Determinant
8. Relations
9. Spaces and Vectors
10. Function and Graphs (Standard)
11. Linear programming
12 Equation in 3-D space
13. Quadratic equation

1. Basic coordinate
In mathematics, the Cartesian coordinate system (also called rectangular coordinate system) is used to determine each point uniquely in a plane through two
numbers, usually called the x-coordinate or abscissa and the y-coordinate or ordinate of the point. To define the coordinates, two perpendicular directed lines (the x-
axis, and the y-axis), are specified, as well as the unit length, which is marked off on the two axes (see Figure 1). Cartesian coordinate systems are also used in
space (where three coordinates are used) and in higher dimensions.

Using the Cartesian coordinate system, geometric shapes (such as curves) can be described by algebraic equations, namely equations satisfied by the coordinates of
the points lying on the shape. For example, the circle of radius 2 may be described by the equation x2 + y2 = 4 (see Figure 2).Contents

Cartesian means relating to the French mathematician and philosopher René Descartes (Latin: Cartesius), who, among other things, worked to merge algebra and
Euclidean geometry. This work was influential in the development of analytic geometry, calculus, and cartography.
The idea of this system was developed in 1637 in two writings by Descartes and independently by Pierre de Fermat, although Fermat did not publish the discovery.[1]
In part two of his Discourse on Method, Descartes introduces the new idea of specifying the position of a point or object on a surface, using two intersecting axes as
measuring guides.[citation needed] In La Géométrie, he further explores the above-mentioned concepts.[2]

Fig. 3 - The four quadrants of a Cartesian coordinate system. The arrows on the axes indicate that they extend forever in their respective directions (i.e. infinitely).
A Cartesian coordinate system in two dimensions is commonly defined by two axes, at right angles to each other, forming a plane (an xy-plane). The horizontal axis is
normally labeled x, and the vertical axis is normally labeled y. In a three dimensional coordinate system, another axis, normally labeled z, is added, providing a third
dimension of space measurement. The axes are commonly defined as mutually orthogonal to each other (each at a right angle to the other). (Early systems allowed
"oblique" axes, that is, axes that did not meet at right angles, and such systems are occasionally used today, although mostly as theoretical exercises.) All the points
in a Cartesian coordinate system taken together form a so-called Cartesian plane. Equations that use the Cartesian coordinate system are called Cartesian
equations.

The point of intersection, where the axes meet, is called the origin normally labeled O. The x and y axes define a plane that is referred to as the xy plane. Given each
axis, choose a unit length, and mark off each unit along the axis, forming a grid. To specify a particular point on a two dimensional coordinate system, indicate the x
unit first (abscissa), followed by the y unit (ordinate) in the form (x,y), an ordered pair.

The choice of letters comes from a convention, to use the latter part of the alphabet to indicate unknown values. In contrast, the first part of the alphabet was used to
designate known values.

An example of a point P on the system is indicated in Figure 3, using the coordinate (3,5).

The intersection of the two axes creates four regions, called quadrants, indicated by the Roman numerals I (+,+), II (−,+), III (−,−), and IV (+,−). Conventionally, the
quadrants are labeled counter-clockwise starting from the upper right ("northeast") quadrant. In the first quadrant, both coordinates are positive, in the second
quadrant x-coordinates are negative and y-coordinates positive, in the third quadrant both coordinates are negative and in the fourth quadrant, x-coordinates are
positive and y-coordinates negative (see table below.)

[edit] Three-dimensional coordinate system

Fig. 4 - Three dimensional Cartesian coordinate system with y-axis pointing away from the observer.

Fig. 5 - Three dimensional Cartesian coordinate system with the x-axis pointing towards the observer.
The coordinate surfaces of the Cartesian coordinates (x, y, z). The z-axis is vertical and the x-axis is highlighted in green. Thus, the red plane shows the points with
x=1, the blue plane shows the points with z=1, and the yellow plane shows the points with y=-1. The three surfaces intersect at the point P (shown as a black sphere)
with the Cartesian coordinates (1.0, -1.0, 1.0).

The three dimensional Cartesian coordinate system provides the three physical dimensions of space — length, width, and height. Figures 4 and 5 show two common
ways of representing it.

The three Cartesian axes defining the system are perpendicular to each other. The relevant coordinates are of the form (x,y,z). As an example, figure 4 shows two
points plotted in a three-dimensional Cartesian coordinate system: P(3,0,5) and Q(−5,−5,7). The axes are depicted in a "world-coordinates" orientation with the z-axis
pointing up.

The x-, y-, and z-coordinates of a point can also be taken as the distances from the yz-plane, xz-plane, and xy-plane respectively. Figure 5 shows the distances of
point P from the planes.

The xy-, yz-, and xz-planes divide the three-dimensional space into eight subdivisions known as octants, similar to the quadrants of 2D space. While conventions
have been established for the labelling of the four quadrants of the x-y plane, only the first octant of three dimensional space is labelled. It contains all of the points
whose x, y, and z coordinates are positive.

The z-coordinate is also called applicate.

Orientation and handedness

The right hand rule.

Fixing or choosing the x-axis determines the y-axis up to direction. Namely, the y-axis is necessarily the perpendicular to the x-axis through the point marked 0 on the
x-axis. But there is a choice of which of the two half lines on the perpendicular to designate as positive and which as negative. Each of these two choices determines
a different orientation (also called handedness) of the Cartesian plane.

The usual way of orienting the axes, with the positive x-axis pointing right and the positive y-axis pointing up (and the x-axis being the "first" and the y-axis the
"second" axis) is considered the positive or standard orientation, also called the right-handed orientation.

A commonly used mnemonic for defining the positive orientation is the right hand rule. Placing a somewhat closed right hand on the plane with the thumb pointing up,
the fingers point from the x-axis to the y-axis, in a positively oriented coordinate system.

The other way of orienting the axes is following the left hand rule, placing the left hand on the plane with the thumb pointing up.

Regardless of the rule used to orient the axes, rotating the coordinate system will preserve the orientation. Switching the role of x and y will reverse the orientation.

In three dimensions
Once the x- and y-axes are specified, they determine the line along which the z-axis should lie, but there are two possible directions on this line. The two possible
coordinate systems which result are called 'right-handed' and 'left-handed'. The standard orientation, where the xy-plane is horizontal and the z-axis points up (and
the x- and the y-axis form a positively oriented two-dimensional coordinate system in the xy-plane if observed from above the xy-plane) is called right-handed or
positive.

The name derives from the right-hand rule. If the index finger of the right hand is pointed forward, the middle finger bent inward at a right angle to it, and the thumb
placed at a right angle to both, the three fingers indicate the relative directions of the x-, y-, and z-axes in a right-handed system. The thumb indicates the x-axis, the
index finger the y-axis and the middle finger the z-axis. Conversely, if the same is done with the left hand, a left-handed system results.

Figure 7 is an attempt at depicting a left- and a right-handed coordinate system. Because a three-dimensional object is represented on the two-dimensional screen,
distortion and ambiguity result. The axis pointing downward (and to the right) is also meant to point towards the observer, whereas the "middle" axis is meant to point
away from the observer. The red circle is parallel to the horizontal xy-plane and indicates rotation from the x-axis to the y-axis (in both cases). Hence the red arrow
passes in front of the z-axis.

Figure 8 is another attempt at depicting a right-handed coordinate system. Again, there is an ambiguity caused by projecting the three-dimensional coordinate system
into the plane. Many observers see Figure 8 as "flipping in and out" between a convex cube and a concave "corner". This corresponds to the two possible orientations
of the coordinate system. Seeing the figure as convex gives a left-handed coordinate system. Thus the "correct" way to view Figure 8 is to imagine the x-axis as
pointing towards the observer and thus seeing a concave corner.

A point in space in a Cartesian coordinate system may also be represented by a vector, which can be thought of as an arrow pointing from the origin of the coordinate
system to the point. If the coordinates represent spatial positions (displacements) it is common to represent the vector from the origin to the point of interest as . In
three dimensions, the vector from the origin to the point with Cartesian coordinates (x,y,z) is sometimes written as[3]:

where , , and are unit vectors that point the same direction as the x, y, and z axes, respectively. This is the quaternion representation of the vector, and was
introduced by Sir William Rowan Hamilton. The unit vectors , , and are called the versors of the coordinate system, and are the vectors of the standard basis in three-
dimensions.

2. Straight Lines

Straight-line equations, or "linear" equations, graph as straight lines, and have simple variables with no exponents on them. If you see an equation with only x and y
(as opposed to, say x2 or sqrt(y), then you're dealing with a straight-line equation.
There are different types of "standard" formats for straight lines; the "standard" format your book refers to may differ from that used in some other books. These
"standard" forms are often holdovers from a few centuries ago, when mathematicians couldn't handle very complicated equations, so they tended to obsess about the
simple cases. Don't worry too much about the "standard" forms.
I think the most useful form of straight-line equations is the "slope-intercept" form:
y = mx + b
This is called the slope-intercept form because "m" is the slope and "b" gives the y-intercept. (For a review of how this works, look at slope and graphing.)
I like slope-intercept form the best. It is in the form "y=", which makes it easiest to plug into, either for graphing or doing word problems. Just plug in your x-value; the
equation is already solved for y. Also, this is the only format you can plug into your (nowadays obligatory) graphing calculator; you have to have a "y=" format to use a
graphing utility. But the best part about the slope-intercept form is that you can read off the slope and the intercept right from the equation. This is great for graphing,
and can be quite useful for word problems. Copyright © Elizabeth Stapel 2000-2007 All Rights Reserved

You will be given problems where they give you some pieces of information about a line, and they want you to come up with the equation of the line. How do you do
that? You plug in whatever they give you, and solve for whatever you need, like this:
• Find the equation of the straight line that has slope m = 4
and passes through the point (–1, –6).
 Okay, they've given me the slope m. In giving me a point, they have also given me an x-value and a y-value: x = –1 and y = –6. In the slope-intercept form
of a straight line, I have y, m, x, and b. So the only thing I don't have a value for is b (which gives me the y-intercept). Then all I need to do is plug in my
slope and the x and y from this particular point, and then solve for b:
y = mx + b
(–6) = (4)(–1) + b
–6 = –4 + b
–2 = b
Then the line equation must be "y = 4x – 2".
What if they don't give you the slope?
• Find the equation of the line that passes through the points (–2, 4) and (1, 2).
Well, if I have two points on a straight line, I can always find the slope; that's what the slope formula is for.

Now I have the slope and two points. I know I can find the equation (by solving first for "b") if I have a point and the slope. So I need to pick one of the
points (it doesn't matter which one), and use it to solve for b. Using the point (–2, 4), I get:
y = mx + b
4 = (– 2/3)(–2) + b
4 = 4/3 + b
4 – 4/3 = b
12/3 – 4/3 = b
b = 8/3
...so y = (– 2/3)x + 8/3.
On the other hand, if I use the point (1, 2), I get:
y = mx + b
2 = (– 2/3)(1) + b
2 = – 2/3 + b
2 + 2/3 = b
6/ + 2/ = b
3 3
b = 8/3
...so y = (– 2/3)x + 8/3, just as when I used the other
point.
As you can see, once you have the slope, it doesn't matter which point you use in order to find the line equation. The answer will work out the same either way.

3. Simultaneous Equations
Simultaneous equations are a set of equations which have more than one value which has to be found. At GCSE, it is unlikely that you will have more than two
equations with 2 values (x and y) which need to be found.

Example:
A man buys 3 fish and 2 chips for £2.80
A woman buys 1 fish and 4 chips for £2.60
How much are the fish and how much are the chips?

There are two methods of solving simultaneous equations. Use the method which you prefer.

Method 1: elimination
First form 2 equations. Let fish be f and chips be c.
We know that:
3f + 2c = 280 (1)
f + 4c = 260 (2)
Doubling (1) gives:
6f + 4c = 560 (3)
(3)-(2) is 5f = 300
\ f = 60
Therefore the price of fish is 60p

Substitute this value into (1):

3(60) + 2c = 280
\ 2c = 100
c = 50
Therefore the price of chips is 50p

Method 2: Substitution
Rearrange one of the original equations to isolate a variable.
Rearranging (2): f = 260 - 4c
Substitute this into the other equation:
3(260 - 4c) + 2c = 280
\ 780 - 12c + 2c = 280
\ 10c = 500
\ c = 50
Substitute this into one of the original equations to get f = 60 .

Harder simultaneous equations:

To solve a pair of equations, one of which contains x², y² or xy, we need to use the method of substitution.

Example:
2xy + y = 10 (1)
x+y=4 (2)
Take the simpler equation and get y = .... or x = ....
from (2), y = 4 - x (3)
this can be substituted in the first equation. Since y = 4 - x, where there is a y in the first equation, it can be replaced by 4 - x .
sub (3) in (1), 2x(4 - x) + (4 - x) = 10
\ 8x - 2x² + 4 - x - 10 = 0
\ 2x² - 7x + 6 = 0
\ (2x - 3)(x - 2) = 0
\ either 2x - 3 = 0 or x - 2 = 0
therefore x = 1.5 or 2 .

Substitute these x values into one of the original equations.

When x = 1.5, y = 2.5
when x = 2, y = 2

Simultaneous equation can also be solved by graphical methods.

4. Polynomial and Algebra of polynomial

By now, you should be familiar with variables and exponents. You may have dealt with expressions like 3x4 or 6x. Polynomials are sums of these expressions. Each
piece of the polynomial, each part that is being added, is called a "term". Polynomial terms have variables to whole-number exponents; there are no square roots of
exponents, no fractional powers, and no variables in the denominator. Here are some examples:
6x –2 NOT a polynomial term This has a negative exponent.
1
/x2 NOT a polynomial term This has the variable in the denominator.

sqrt(x) NOT a polynomial term This has the variable inside a radical.

4x2 a polynomial term

Here is a typical polynomial:

Notice the exponents on the terms. The first term has exponent 2; the second term has an understood exponent 1; and the last term doesn't have any variable at all.
Polynomials are usually written this way, with the terms written in "decreasing" order; that is, with the highest exponent first, the next highest next, and so forth, until
you get down to the plain old number.
Any term that doesn't have a variable in it is called a "constant" term because, no matter what value you may put in for the variable
change. In the picture above, no matter what x might be, 7 will always be just 7.
The first term in the polynomial, when it is written in decreasing order, is also the term with the biggest exponent, and is called the "leading term".
The exponent on a term tells you the "degree" of the term. For instance, the leading term in the above polynomial is a "second-degree term". The second term is a
"first degree" term. The degree of the leading term tells you the degree of the whole polynomial; the polynomial above is a "second-degree polynomial". Here are a
couple more examples:
• Give the degree of the following polynomial: 2x5 – 5x3 – 10x + 9
This polynomial has four terms, including a fifth-degree term, a third-degree term, a first-degree term, and a constant term.
This is a fifth-degree polynomial.
• Give the degree of the following polynomial: 7x4 + 6x2 + x
This polynomial has three terms, including a fourth-degree term, a second-degree term, and a first-degree term. There is no constant term.
This is a fourth-degree polynomial.

When a term contains both a number and a variable part, the number part is called the "coefficient". The coefficient on the leading term is called the leading
coefficient.
In the above example, the coefficient of the leading term is 4; the coefficient of the second term is 3; the constant term doesn't have a coefficient.
Elizabeth Stapel 2000-2007 All Rights Reserved
The "poly" in "polynomial" means "many". I suppose, technically, the term "polynomial" should only refer to sums of many terms, but the term is used to refer to
anything from one term to a zillion terms. However, the shorter polynomials do have their own names:
• a one-term polynomial, such as 2x or 4x2, may also be called a "monomial"

• a two-term polynomial, such as 2x + y or x2 – 4, may also be called a "binomial"

• a three-term polynomial, such as 2x + y + z or x4 + 4x2 – 4, may also be called a "trinomial"
I don't know if there are names for polynomials with a greater numbers of terms; I've never heard of any names other than what I've listed.

Polynomials are also sometimes named for their degree:

• a second-degree polynomial, such as 4x2, x2 – 9, or ax2 + bx + c, is also called a "quadratic"
• a third-degree polynomial, such as –6x3 or x3 – 27, is also called a "cubic"
• a fourth-degree polynomial, such as x4 or 2x4 – 3x2 + 9, is sometimes called a "quartic"
• a fifth-degree polynomial, such as 2x5 or x5 – 4x3 – x + 7, is sometimes called a "quintic"
There are names for some of the higher degrees, but I've never heard of any names being used other than the ones I've listed.

"Evaluating" a polynomial is the same as evaluating anything else: you plug in the given value of x, and figure out what y is supposed to be. For instance:
• Evaluate 2x3 – x2 – 4x + 2 at x = –3
Plug in –3 for x, remembering to be careful with parentheses and negatives:
2(–3)3 – (–3)2 – 4(–3) + 2
= 2(–27) – (9) + 12 + 2
= –54 – 9 + 14
= –63 + 14
= –49

6. Matrices and Determinant

A determinant is a square array of elements with a prescribed evaluation rule which is depicted as follows

The evaluation rule can be illustrated by a determinant of order three.

Determinants are useful in the evaluation of vector products and vector operations like the curl.

Determinant Evaluation Example

For a determinant of order three the evaluation rule is

Take the elements of the top row and multiply them times the determinant of their cofactors. The cofactor is the array left when the row and column of the given top
row element is eliminated. The evaluation of the determinant of the cofactor follows the same pattern until the cofactor has dimension two. At that point, it's value is
the difference of the diagonal products.
With each square matrix corresponds just one number. This number is called the determinant of the matrix. The determinant of a matrix A is denoted det(A) or |A|.
Now we'll define this correspondence.
Determinant of a 1 x 1 matrix
De determinant of the matrix is the element itself.
Ex: det([-7]) = -7
Permutation of n ordered elements
Say S is an ordered set of n elements. A one-one transformation t of the set S onto itself is a permutation of S.
Example: S = (1, 2, 3, 4, 5) . A permutation t is defined by
t(1, 2, 3, 4, 5) = (2, 5, 4, 1, 3)
The permutation, witch transforms (2, 5, 4, 1, 3) back to (1, 2, 3, 4, 5) is called the inverse permutation of t.
Transposition
A permutation, which interchanges two elements, and fixes all others, is called a transposition.
Example: S = (1, 2, 3, 4, 5) . The permutation defined by
t(1, 2, 3, 4, 5) = (1, 4, 3, 2, 5) is a transposition
Theorem
Every permutation of n ordered elements can be expressed as a sequence of transpositions. If this permutation is a sequence of an even number of transpositions, it
is impossible to write this permutation as a sequence of an odd number of transpositions.
Even and odd permutations
If a permutation of n ordered elements can be expressed as an even number of transpositions, then it is called an even permutation. If a permutation of n ordered
elements can be expressed as an odd number of transpositions, then it is called an odd permutation.
The inverse permutation has exactly the same number of transpositions as t.
So, if t is even, its inverse is even too.
Example: S = (1, 2, 3, 4, 5)
t(1, 2, 3, 4, 5) = (1, 3, 4, 2, 5) is an even permutation.
t(1, 2, 3, 4, 5) = (1, 3, 4, 5, 2) is an odd permutation.
Sign of a permutation
The sign of an even permutation t is +1. We write : sgn(t) = +1.
The sign of an odd permutation t is -1. We write : sgn(t) = -1.
The determinant of an n x n matrix
Let S = (1, 2, 3, ... , n)
t is a permutation of S, so t(1, 2, 3, ... , n) = (t(1), t(2), ... , t(n)).
There are n! permutations of S.
A is a n x n matrix with elements ai, j.

Now, with each permutation t of S, create the product

sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) . ... . an, t(n). There are n! such products.
|A| is defined as the sum of all those products.
Note that each term of |A| involves each row and each column only once.
Example1 : We want to calculate the determinant of a 2x2 matrix A.
Now n = 2 and there are only two permutations of S = (1, 2).
t(1, 2) = (1, 2) with sgn(t) = +1
t'(1, 2) = (2, 1) with sign(t') = -1
We have only two terms +1.a1, 1 . a2, 2 and -1.a1, 2 . a2, 1
Thus the determinant of A is a1, 1 . a2, 2 - a1, 2 . a2, 1
We don't forget the rule :

|a b|
|c d|

= ad - cb
Example2 : We want to calculate the determinant of a 3x3 matrix A.
Now n = 3 and there are only 6 permutations of S = (1, 2, 3).
These 6 permutations transform (1, 2, 3) in:

(1, 2, 3) (2, 3, 1) (3, 1, 2) (even permutations)

(3, 2, 1) (1, 3, 2) (2, 1, 3) (odd permutations)
Now we have six terms to add
a1, 1 . a2, 2 . a3, 3 + a1, 2 . a2, 3 . a3, 1 + a1, 3 . a2, 1 . a3, 2
-a1, 3 . a2, 2 . a3, 1 - a1, 1 . a2, 3 . a3, 2 - a1, 2 . a2, 1 . a3, 3
We don't forget the rule :

|a b c|
|d e f|
|g h i|

= aei + bfg + cdh - ceg - afh - bdi

The last rule is known as the Sarrus rule for 3 x 3 determinants.
To calculate larger determinants there are a lot of other methods involving various properties of determinants.
Row and columns of the determinant
If we say the ith row of a determinant we mean the ith row of the matrix corresponding with this determinant. If we say the ith column of a determinant we mean the
ith column of the matrix corresponding with this determinant.
Cofactor of an element ai, j
Now choose a fixed row value i.
Since each row appears once and only once in each term of |A|, each term of |A| contains exactly one of the factors ai, 1, ai, 2, ai, 3
Thus, we can write |A| as a linear polynomial in ai, 1, ai, 2, ai, 3, ... ai, n.
We denote the coefficients respectively Ai, 1, Ai, 2, Ai, 3, ... Ai, n.
These coefficients are called the cofactors. Ai, j is the cofactor of ai, j.
|A| = Ai, 1 . ai, 1 + Ai, 2 . ai, 2 + Ai, 3 . ai, 3 + ... Ai, n . ai, n.
Since each term of |A| involves each row and each column only once, the cofactor Ai, j is independent of the elements of the ith row and the elements of the jth
column.
It contains only elements from the matrix obtained from A by crossing out the ith row and the jth column.
Remark: If we write |A| = Ai, 1 . ai, 1 + Ai, 2 . ai, 2 + Ai, 3 . ai, 3 + ... Ai, n . ai, n, we say that the determinant is calculated emanating from the ith row.
Example :

|a b c|
|d e f|
|g h i|

= aei + bfg + cdh - ceg - afh - bdi

Choose for instance row 2.
Each term of |A| contains exactly one of the factors d , e ,f .
Thus, we can write |A| as a linear polynomial in d , e, f.
|A| = (ch-bi)d + (ai-cg)e + (bg-ah)f
ch-bi is the cofactor of d.
ai-cg is the cofactor of e.
bg-ah is the cofactor of f.
Not any cofactor contains an element of the chosen row 2.
The cofactor of d contains neither an element of row 2 nor an element of column 1.

If we write |A| = (ch-bi)d + (ai-cg)e + (bg-ah)f , we say that the determinant is calculated emanating from the second row.
Similarly, we can start with a fixed column and then write |A| as a linear polynomial in a1, j, a2, j, a3, j, ... an, j. Then one finds the same cofactors. So a
cofactor Ai, j.
A matrix A and its transpose have the same determinant.
We call A' the transpose of A, ai, j' = aj, i.
We know that |A'|
= sum of all products sgn(t) . a1, t(1)' . a2, t(2)' . a3, t(3)' ... an, t(n)'
= sum of all products sgn(t)at(1), 1 . at(2), 2 . at(3), 3 ... at(n), n.
Since t(1), t(2), ... , t(n) is a permutation of 1, 2, 3 ... n , we can reorder the factors of each term, according to the first index. This can be done using the inverse
permutation of t. The permutation t transforms (1, 2, 3 ... n) to (t(1), t(2), ... , t(n)), so the inverse permutation t' brings (t(1), t(2), ... , t(n)) back to (1, 2, 3 ... n) and this
inverse permutation has exactly the same number of transpositions as t. So sign(t) = sign(t'). Then |A'|
= sum of all products sgn(t') . a1, t'(1) . a2, t'(2) . a3, t'(3) ... an, t'(n)
Because the set of all permutations is the same set of all inverse permutations. |A| = |A'|.
Important result
Appealing on previous property, it is immediate that each property we'll find for the rows of a matrix, also holds for the columns and each property for the columns
holds for the rows.
Interchanging two columns of A
First, denote t' the permutation transposing only i and j.
Thus t'(1, ... ,i, ... , j, ... , n) = (1, ... , j, ... , i, ... , n). Sgn(t') = -1 and for each permutation t we have sign(t't) = -sign(t).
Say A' is obtained by interchanging the column i and j of A.
For each k we have ak, i' = ak, j or even for each k and each l we have ak, l' = ak, t'(l)
We investigate |A'|.
We know that |A'|
= sum of all products sgn(t) . a1, t(1)' . a2, t(2)' . a3, t(3)' ... an, t(n)'
= sum of all products sgn(t) . a1, t't(1) . a2, t't(2) . a3, t't(3) ... an, t't(n)
= sum of all products -sgn(t't) . a1, t't(1) . a2, t't(2) . a3, t't(3) ... an, t't(n)
Since the set of permutations of (1 ... n) is a group, the set of all permutations t and the set of all permutations t" = t't is the same set.
Therefore |A'|
= sum of all products -sgn(t't) . a1, t't(1) . a2, t't(2) . a3, t't(3) ... an, t't(n)
= sum of all products -sgn(t") . a1, t"(1) . a2, t"(2) . a3, t"(3) ... an, t"(n)
= -|A|
Conclusion :
When we change two columns in A, |A| changes sign.
When we change two rows in A, |A| changes sign.
Multiplying a row of A with a real number
Say A' is obtained by multiplying the ith row of A by a real number r.
Then ai, k' = r . ai, k for each k and fixed i, and if j is not i then aj, k' = aj, k
We know that |A'|
= sum of all products sgn(t) . a1, t(1)' . a2, t(2)' . a3, t(3)' ... ai, t(i)' ... an, t(n)'
= sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... r.ai, t(i) ... an, t(n)
= sum of all products r . sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... ai, t(i) ... an, t(n)
= r.(sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... ai, t(i) ... an, t(n))
= r.|A|
Conclusion :
When we multiply a row in A with a real number r, |A| changes in r.|A|
When we multiply a column in A with a real number r, |A| changes in r.|A|
If A has two equal rows, |A| = 0
If we interchange these two rows, the determinant does not change. Appealing on previous property the determinant changes in its opposite. This is only possible
when the determinant = 0.
Addition of two determinants, which differ only from the ith row
Say A and B are matrices which are only different in the ith row.
For all j different from i and all k we have aj,k = bj, k.
|A|
= sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... ai, t(i) ... an, t(n)
|B|
= sum of all products sgn(t) . b1, t(1) . b2, t(2) . b3, t(3) ... bi, t(i) ... bn, t(n)
= sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... bi, t(i) ... an, t(n)
So |A| + |B|
= sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... (ai, t(i) + bi, t(i)) ... an, t(n)
= determinant of the matrix formed by adding the ith row from A and B, and taking the other elements from A or from B.
The same rule holds for columns
Ex.

|a b c| |a b' c| |a b+b' c|
|d e f|+|d e' f| = |d e+e' f|
|g h i| |g h' i| |g h+h' i|
Equal determinants
Let A be any square matrix.
Let B be the matrix formed by replacing in A the ith row with the jth row, leaving the jth row unchanged.
Since B has two equal rows |B| = 0.
Let C be the matrix formed by multiplying the ith row from B with r, leaving the other elements unchanged.
Then |C| = r.|B|=0
A an C differ only from the ith row, so we can use previous property and we have:
|A|+0 = |A|+|C| = determinant of the matrix formed by adding the ith row from A and C, and taking the other elements from A or from C.
Therefore, a determinant does not change if we add a multiple of a row to another row. The same rule holds for columns
Ex.

|a b c| |a+rd b+re c+rf|

|d e f| = | d e f |
|g h i| | g h i |
The cofactor A1, 1
Let A be any square matrix. We know that |A|
= sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... an, t(n).
A1, 1 is the coefficient of a1, 1 in this sum.
The terms containing a1, 1 are the terms with t(1) = 1.
Recall the sum of all products sgn(t) . a1, t(1) . a2, t(2) . a3, t(3) ... an, t(n).
Instead of taking this sum over all permutations of (1 ... n), we only take the sum over the permutations t' with t(1) = 1.
This sum then gives A1, 1 . a1, 1.
These special permutations t' are all the permutations of (2..n).
Thus, A1, 1 . a1, 1= sum of all products sgn(t') . a1, 1 . a2, t'(2) . a3, t'(3') ... an, t'(n).
Then A1, 1 = sum of all products sgn(t') a2, t'(2) . a3, t'(3') ... an, t'(n).
This is, by the definition of a determinant, the determinant of the sub-matrix of A obtained from A by crossing out the first row and the first column.
Conclusion:
A1, 1 = the determinant of the sub-matrix of A obtained from A by crossing out the first row and the first column.
The cofactor Ai, j
Let A be any square matrix. Focus the element e = ai, j. Interchange in succession row i and i-1; i-1 and i-2; ... until e is on the first row. This demands i-1 steps. Then
we interchange in succession column j and j-1; j-1 and j-2; ... until e is on the first column and on the first row. This demands j-1 steps. During this process the
determinant of the matrix changes i+j-2 times sign. Now the cofactor of e is the determinant of the sub-matrix of obtained from by crossing out the first row and the
first column. Now return to the original matrix.
The value of Ai, j= (-1)i+j-2.(the determinant of the sub-matrix of A, obtained from A by crossing out the ith row and the jth column.
Or stated simpler: The value of Ai, j = (-1)i+j.(the determinant of the sub-matrix of A, obtained from A by crossing out the ith row and the jth column.
|I| = 1

8. Relations
A relation R on A is defined to be R A x A
• reflexive

For all a A (a,a) R

• A reflexive relation R on set X is one where for all a in X, a is R-related to itself.
• Examples include equality =, < greater/less than or equal to, “is a subset of” , divides |
• An irreflexive relation R is one where for all a in X, a is never R-related to itself
• anti-symmetric

• A binary relation R on a set X is antisymmetric if, for all a and b in X, if a is related to b and b is related to a, then a = b
• To prove antisymmetric show a1< a2,a2< a1 so a1=a2
Note that 'antisymmetric' is not the logical negative of 'symmetric' (whereby aRb implies bRa).
• A binary relation R over a set X is symmetric if it holds for all a and b in X that if a is related to b then b is related to a.
• A bijection A-> B is symmetric since inverse of bijection is a bijection
• symmetric
• transitive
• A binary relation R over a set X is transitive if it holds for all a, b, and c in X, that if a is related to b and b is related to c, then a is related to c.
• eg divisibiltiy is transitive: a|b and b|c => a|c
• A bijection A-> B is transitive because composition of bijections is a bijection

Partial order:
A partial order is a binary relation R over a set P which is reflexive, antisymmetric, and transitive

Examples:
• The set of natural numbers equipped with the lesser than or equal to relation, or with the divisibility relation is a partial order. Under divisibility there is a
lower (lcm) and upperbound (gcd)
• If (A, <A) and (B,<B) are partially ordered sets, then the product order is
(a1 <A a2) ^ (b1 <A b2) <-> (a1, b1) < (a2, b2)
ie Its just the Cartesian (normal multiplication) product of the two orders.
You can prove AxB is a partial order as follows:
• reflexive as (a,b) < (a,b)
• A and B are antisymmetric, so a1< a2 and a2< a1 so a1=a2 (same with b's) so
(a1,b1)=(a2,b2)
• a1<Aa2<Aa3 so a1<Aa3 because <A is transitive, same with b's
(a1,b1)

Divisibility is a partial order:

• reflexive a|a
• antisymmetric a|b => b|a
• transitive a|b and b|c => a|c

To prove something is a least upper bound, assume there is a lower upper bound, then it divides the elements and either divides the LUB or is the LUB:
contradiction. Simliar argument for greatest lower bound.
is a partial order:
• reflexive AA
• anti symmetric A B B A=> B=A
• transitive A B B C => A C
Given A,B X the greatest lower bound is A B and the lowest upper bound is A n B.
This is proved by letting C be g|b, C X and by entailment show either CAB or A n BC
The greatest lower bound of AxB = ( glb(a1,a2) ,glb(b1,b2) )
Subsets including of 0 dont have an upper bound, as 0 doesn't divide natural numbers (when using the divide operator as the PO).
N has no least upper bound (as there is no number bigger than the largest in N), but a greatest lower bound. Its subsets have both.
Lattices:
A lattice is a partially ordered set where every pair of elements has both a least upper bound and a greatest lower bound.

Well founded relation

• A well founded relation on a set A is a relation A x A such that there are no infinitely descending chains.
• Any non-empty subset S of a well-founded (relation) set A has a minimal element: Pick an element a0, if it is minimal then stop. Otherwise pick a1 a0. If
minimal stop, if not then pick a2 a1. This can't continue for ever as is well founded, therefore there must be a minimal element.
• If two relations <A on A and <B on B are well founded show that the relation and product relation on AxB are both well founded:
Consider projection of product order into A -> find minimum.
Consider proection into B-> find minimum.
Show minimal
Observe <L <P
Total order
A total order on a set X is any binary relation on X that is antisymmetric, transitive, and total. This means that if we denote one such relation by = then the following
statements hold for all a, b and c in X:
• if a = b and b = a then a = b (antisymmetry)
• if a = b and b = c then a = c (transitivity)
• a = b or b = a (totality)
A relation's property of "totality" can be described this way: that any pair of elements in the set are mutually comparable under the relation.
Notice that the totality condition implies reflexivity, that is a = a. Thus a total order is also a partial order, that is, a binary relation which is reflexive, antisymmetric and
transitive.

• The lexicographical order on the product of a set of totally ordered sets indexed by an ordinal (indexing number, eg first, second etc.) is itself a total
order.
• For example, any set of words ordered alphabetically is a total order, viewed as a subset of a product of a countable number of copies of a set formed by
adding the space symbol to the alphabet (and defining a space to be less than any letter).
• A totally ordered set is said to be complete if every subset that has an upper bound, has a least upper bound.
Equivalence relation
An equivalence relation on a set X is a binary relation on X that is reflexive, symmetric, and transitive. That is, if the relation is denoted by the symbol "~", it holds for
all a, b, and c in X that
1. (Reflexivity) a ~ a
2. (Symmetry) if a ~ b then b ~ a
3. (Transitivity) if a ~ b and b ~ c then a ~ c
• If two equivalence relations over the set X are given, then their intersection (viewed as subsets of X×X) is also an equivalence relation.

9. Spaces and Vectors

Vector spaces and systems of linear equations

Vector spaces and homogeneous systems

Take a homogeneous system of linear equations in n unknowns. Each solution of that system can be viewed as a vector from the vector space V of all the real n-
tuples.
Each real multiple of that solution is a solution too, and the sum of two solutions is a solution too. Therefore, all the solutions of the system form a subspace M of V. It
is called the solution space of the system.

Basis of a solution space

By means of an example, we show how a basis of a solution space can be found.

/ 2x + 3y - z + t = 0
\ x - y + 2z - t = 0

This is a system of the second kind.

x and y can be taken as main unknowns.
z and t are the side unknowns.
The solutions are

x = -z + (2/5)t
y = z - (3/5)t
The set of solutions can we written as
(-z + (2/5)t , z - (3/5)t , z , t ) with z and t in R

<=>

z(-1,1,1,0) + t(2/5,-3/5,0,1) with z and t in R

Hence, all solutions are linear combinations of the linear independent vectors (-1,1,1,0) and (2/5,-3/5,0,1).
These two vectors constitute a basis of the solution space.

Solutions of a non homogeneous system

We can denote such system shortly as AX = B, with coefficient matrix A, the column matrix B of the known terms and X is the column matrix of the unknowns.
Consider also the corresponding homogeneous system AX = 0 with the same A and X as above.
If X' is a fixed solution of AX = B then AX' = B .
If X" is a arbitrary solution of AX = 0 then AX" = 0 .
Then, AX' + AX" = B <=> A(X' + X") = B <=> X' + X" is a solution of AX = B.
Conclusion: If we add an arbitrary solution of AX = 0 to a fixed solution of AX = B then X' + X" is a solution of AX = B.
Furthermore:
If X' is a fixed solution of AX = B then AX' = B .
If X" is a arbitrary solution of AX = B then AX" = B .
Then, AX" - AX' = 0 <=> A(X" - X') = 0 <=> X" - X' is a solution of AX = 0.
<=> X" = X' + (a solution of AX = 0).
Conclusion:
Any arbitrary solution of AX = B is the sum of a fixed solution of AX = B and a solution of AX = 0
So, if we have a fixed solution of AX = B and we add to this solution all the solutions of the corresponding homogeneous system one after another, then we get all
solutions AX = B.
Example:

/ 2x + 3y - z + t = 0
\ x - y + 2z - t = 0

Above we have seen that the solutions are z(-1,1,1,0) + t(2/5,-3/5,0,1).

/ 2x + 3y - z + t = 5
\ x - y + 2z - t = 0

has a solution (1,1,0,0) .

All solutions of the last system are (1,1,0,0) + z(-1,1,1,0) + t(2/5,-3/5,0,1).

Sum of two vector spaces

Say A and B are subspaces of a vector space V.
We define the sum of A and B as the set
{ a + b with a in A and b in B }
We write this sum as A + B.

The sum as subspace

Theorem: The sum A+B, as defined above, is a subspace of V.
Proof:
For all a1 and a2 in A and all b1 and b2 in B and all r, s in R we have

r(a1 + b1) + s(a2 + b2) = (r a1 + s a2) + (r b1 + s b2)

is in A + B.

Direct sum of two vector spaces

The sum A+B, as defined above, is a direct sum if and only if the vector 0 is the only vector common to A and B.

Example

In the space R3
A = span{ (3,2,1) }
B = span{ (2,1,4) ; (0,1,3) }
Investigate if A+B is a direct sum

Say r,s,t are real numbers, then

each vector in space A is of the form r.(3,2,1) and
each vector in space B is of the form s.(2,1,4) + t.(0,1,3) .

For each common vector, there is a suitable r,s,t such that

r.(3,2,1) = s.(2,1,4) + t.(0,1,3)

<=>

r.(3,2,1) - s.(2,1,4) - t.(0,1,3) = (0,0,0)

<=>
 / 3r - 2s = 0
| 2r - s - t = 0
\ r - 4s -3t = 0

Since |3 -2 0|
|2 -1 -1| is not 0,
|1 -4 -3|

the previous system has only the solution r = s = t = 0.

The vector (0,0,0) is the only common vector of A and B.

Thus, A+B is a direct sum.

Property
If A + B is a direct sum, then each vector v in A+B can be written, in just one way, as the sum of an element of A and an element of B.
Proof:

Suppose v = a1 + b1 = a2 + b2 with ai in A and bi in B.

Then a1 - a2 = b2 - b1 and a1 - a2 is in A and b2 - b1 is in B

Therefore a1 - a2 = b2 - b1 is a common element of A and B.

But the only common element is 0.

So, a1 - a2 = 0 and b2 - b1 = 0

and a1 = a2 and b2 = b1

Supplementary vector spaces

Say that vector space V is the direct sum of A and B, then

A is the supplementary vector space of B with respect to V.

B is the supplementary vector space of A with respect to V.

A and B are supplementary vector spaces with respect to V.

Basis and direct sum

Theorem:
Say V is the direct sum of the spaces M and N.
If {a,b,c,..,l } is a basis of M and {a',b',c',..,l' } is a basis of N,
then {a,b,c,..,l,a',b',c',..,l' } is a basis of M+N.

Proof:
Each vector v of V can be written as m + n with m in M and n in N.
Then m = ra + sb + tc + ... + zl and n = r'a' + s'b' + t'c' + ... + z'l' , with r,s,t,...z,r',s',t',...z' real coefficients.
Thus each vector v = ra + sb + tc + ... + zl + r'a' + s'b' + t'c' + ... + z'l'
Therefore the set {a,b,c,..,l,a',b',c',..,l'} generates V.
If ra + sb + tc + ... + zl + r'a' + s'b' + t'c' + ... + z'l' = 0 , then ra + sb + tc + ... + zl is a vector m of M and
r'a' + s'b' + t'c' + ... + z'l' is a vector n of N.
From a previous theorem we know that we can write the vector 0 in just one way, as the sum of an element of M and an element of N. That way is 0 = 0 + 0 with 0 in
M and 0 in N.
From this we see that necessarily m = 0 and n = 0 and thus ra + sb + tc + ... + zl = 0 and
r'a' + s'b' + t'c' + ... + z'l' = 0
Since all vectors in these expressions are linear independent, it is necessarily that all coefficients are 0 and from this we know that the generating vectors
{a,b,c,..,l,a',b',c',..,l'} are linear independent.
Dimension of a direct sum
From previous theorem it follows that
dim(A+B) = dim(A) + dim(B)

Converse theorem
If {a,b,c,..,l } is a basis of M and {a',b',c',..,l' } is a basis of N, and {a,b,c,..,l,a',b',c',..,l' } are linear independent,
then M+N is a direct sum.
Proof:
Each element m of M can be written as ra + sb + tc + ... + zl .
Each element n of N can be written as r'a' + s'b' + t'c' + ... + z'l' .
For a common element we have

ra + sb + tc + ... + zl = r'a' + s'b' + t'c' + ... + z'l'

<=>
ra + sb + tc + ... + zl - r'a' - s'b' - t'c' - ... - z'l' = 0

Since all vectors are linear independent, all coefficients must be 0. The only common vector is 0.

Direct sum criterion

From the two previous theorems we deduce that: If {a,b,c,..,l } is a basis of M and {a',b',c',..,l' } is a basis of N, then

M+N is a direct sum <=> {a,b,c,..,l,a',b',c',..,l' } are linear independent

Projection in a vector space

Choose two supplementary subspaces M and N with respect to the space V. Each vector v of V can be written in exactly one way as the sum of an element m of M
and an element n of N.
Then v = m + n .
Now we can define the transformation

p: V --> V : v --> m

We define this transformation as

the projection of V on M with respect to N

Projection, example
V is the space of all polynomials with a degree not greater than 3.
We define two supplementary subspaces
M = span { 1, x }
N = span { x2, x3 }
Each vector of V is the sum of exactly one vector of M and of N.
e.g. 2x3 - x2 + 4x - 7 = (2x3 - x2) + (4x - 7)
Say p is the projection of V on M with respect to N, then

p(2x3 - x2 + 4x - 7 ) = 4x - 7

Say q is the projection of V on N with respect to M, then

q(2x3 - x2 + 4x - 7 ) = 2x3 - x2

Similarity transformation of a vector space

Let r = any constant real number.
In a vector space V we define the transformation
 h : V --> V : v --> r.v

We say that h is a similarity transformation of V with factor r.

Important special values of r are 0, 1 and -1.

Reflection in a vector space

Choose two supplementary subspaces M and N with respect to the space V.
Each vector v of V is the sum of exactly one vector m of M and n of N.
Now we define the transformation

s : V --> V : v --> m - n

We say that s is the reflection of V in M with respect to the N.

This definition is a generalisation of the ordinary reflection in a plane. Indeed, if you take the ordinary vectors in a plane and if M and N are one dimensional
supplementary subspaces, then you'll see that with the previous definition, s becomes the ordinary reflection in M with respect to the direction given by N.

Example of a reflection

Take V = R4.

M = span{(0,1,3,1);(1,0,-1,0)}

N = span{(0,0,0,1);(3,2,1,0)}

It is easy to show that M and N have only the vector 0 in common. (This is left as an exercise.) So, M and N are supplementary subspaces.
Now we'll calculate the image of the reflection of vector v = (4,3,3,1) in M with respect to N.
First we write v as the sum of exactly one vector m of M and n of N.

(4,3,3,1) = x.(0,1,3,1) + y.(1,0,-1,0) + z.(0,0,0,1) + t.(3,2,1,0)

The solution of this system gives x = 1; y = 1; z = 0; t = 1. The unique representation of v is

(4,3,3,1) = (1,1,2,1) + (3,2,1,0)

The image of the reflection of vector v = (4,3,3,1) in M with respect to N is vector v' =

(1,1,2,1) - (3,2,1,0) = (-2,-1,1,1)

10. Function and Graphs (Standard)

The graph of a function is the set of all points whose co-ordinates (x, y) satisfy the function y = f(x). This means that for each
corresponding y-value which is obtained when we substitute into the expression for f(x).
Since there is no limit to the possible number of points for the graph of the function, we will follow this procedure at first:
- select a few values of x
- obtain the corresponding values of the function
- plot these points by joining them with a smooth curve
However, you are encouraged to learn the general shapes of certain common curves (like straight line, parabola, trigonometric and exponential curves) - it's much
easier than plotting points and more useful for later!

Example 1
A man who is 2 m tall throws a ball straight up and its height at time t (in s) is given by h = 2 + 9t - 4.9t2 m. Graph the function.
Answer
We start at t = 0 since negative values of time have no practical meaning here.
x 0 0.5 1 1.5 2 y 2 5.3 6.1 4.5 0.4
This shape is called a parabola and is common in applications of mathematics.
NOTE:
We could have written the function in this example as: h(t) = 2 + 9t - 4.9t2.

Example 2
The velocity (in m/s) of the ball in Example 1 at time t (in s) is given by
v = 9 - 9.8t
Draw the v-t graph. What is the velocity when the ball hits the ground?
Answer
Since we recognise it is a straight line, we only need to plot 2 points and join them. But we find 3 ponts, just to check.
x 0 1 2 y 9 -0.8 -10.6

The ball hits the ground at approx t = 2.05 s (we can see this from Example 1). The velocity when the ball hits the ground from the graph we just drew is about -11
m/s.
Normally, we take velocity in the up direction to be positive.

11. Linear programming

Linear programming is often a favorite topic for both professors and students. The ability to introduce LP using a graphical approach, the relative ease of the solution
method, the widespread availability of LP software packages, and the wide range of applications make LP accessible even to students with relatively weak
mathematical backgrounds. Additionally, LP provides an excellent opportunity to introduce the idea of "what-if" analysis, due to the powerful tools for post-optimality
analysis developed for the LP model.
Linear Programming (LP) is a mathematical procedure for determining optimal allocation of scarce resources. LP is a procedure that has found practical application
in almost all facets of business, from advertising to production planning. Transportation, distribution, and aggregate production planning problems are the most
typical objects of LP analysis. In the petroleum industry, for example a data processing manager at a large oil company recently estimated that from 5 to 10 percent
of the firm's computer time was devoted to the processing of LP and LP-like models.
Linear programming deals with a class of programming problems where both the objective function to be optimized is linear and all relations among the variables
corresponding to resources are linear. This problem was first formulated and solved in the late 1940's. Rarely has a new mathematical technique found such a wide
range of practical business, commerce, and industrial applications and simultaneously received so thorough a theoretical development, in such a short period of time.
Today, this theory is being successfully applied to problems of capital budgeting, design of diets, conservation of resources, games of strategy, economic growth
prediction, and transportation systems. In very recent times, linear programming theory has also helped resolve and unify many outstanding applications.
It is important for the reader to appreciate, at the outset, that the "programming" in Linear Programming is of a different flavor than the "programming" in Computer
Programming. In the former case, it means to plan and organize as in "Get with the program!", it programs you by its solution. While in the latter case, it means to
write codes for performing calculations. Training in one kind of programming has very little direct relevance to the other. In fact, the term "linear programming" was
coined before the word "programming" became closely associated with computer software. This confusion is sometimes avoided by using the term linear optimization
as a synonym for linear programming.
Any LP problem consists of an objective function and a set of constraints. In most cases, constraints come from the environment in which you work to achieve your
objective. When you want to achieve the desirable objective, you will realize that the environment is setting some constraints (i.e., the difficulties, restrictions) in
fulfilling your desire or objective. This is why religions such as Buddhism, among others, prescribe living an abstemious life. No desire, no pain
advice with respect to your business objective?
What is a function: A function is a thing that does something. For example, a coffee grinding machine is a function that transform the coffee beans into powder. The
(objective) function maps and translates the input domain (called the feasible region) into output range, with the two end-values called the maximum and the
minimum values.
When you formulate a decision-making problem as a linear program, you must check the following conditions:
1. The objective function must be linear. That is, check if all variables have power of 1 and they are added or subtracted (not divided or
multiplied)
2. The objective must be either maximization or minimization of a linear function. The objective must represent the goal of the decision-maker
3. The constraints must also be linear. Moreover, the constraint must be of the following forms ( , , or =, that is, the LP-constraints are always
closed).
For example, the following problem is not an LP: Max X, subject to X  1. This very simple problem has no solution.
As always, one must be careful in categorizing an optimization problem as an LP problem. Here is a question for you. Is the following problem an LP problem?
Max X2
subject to:
X1 + X2  0
X12 - 4  0
Although the second constraint looks "as if" it is a nonlinear constraint, this constraint can equivalently be written as:
X1  -2, and X2  2.
Therefore, the above problem is indeed an LP problem.
For most LP problems one can think of two important classes of objects: The first is limited resources such as land, plant capacity, or sales force size; the second, is
activities such as "produce low carbon steel", "produce stainless steel", and "produce high carbon steel". Each activity consumes or possibly contributes additional
amounts of the resources. There must be an objective function, i.e. a way to tell bad from good, from an even better decision. The problem is to determine the best
combination of activity levels, which do not use more resources than are actually available. Many managers are faced with this task everyday. Fortunately, when a
well-formulated model is input, linear programming software helps to determine the best combination.
The Simplex method is a widely used solution algorithm for solving linear programs. An algorithm is a series of steps that will accomplish a certain task.

LP Problem Formulation Process and Its Applications

To formulate an LP problem, I recommend using the following guidelines after reading the problem statement carefully a few times.
Any linear program consists of four parts: a set of decision variables, the parameters, the objective function, and a set of constraints. In formulating a given decision
problem in mathematical form, you should practice understanding the problem (i.e., formulating a mental model) by carefully reading and re-reading the problem
statement. While trying to understand the problem, ask yourself the following general questions:
1. What are the decision variables? That is, what are controllable inputs? Define the decision variables precisely, using descriptive names.
Remember that the controllable inputs are also known as controllable activities, decision variables, and decision activities.
2. What are the parameters? That is, what are the uncontrollable inputs? These are usually the given constant numerical values. Define the
parameters precisely, using descriptive names.
3. What is the objective? What is the objective function? Also, what does the owner of the problem want? How the objective is related to his
decision variables? Is it a maximization or minimization problem? The objective represents the goal of the decision-maker.
4. What are the constraints? That is, what requirements must be met? Should I use inequality or equality type of constraint? What are the
connections among variables? Write them out in words before putting them in mathematical form.
Learn that the feasible region has nothing or little to do with the objective function (min or max). These two parts in any LP formulation come mostly from two distinct
and different sources. The objective function is set up to fulfill the decision-maker's desire (objective), whereas the constraints which shape the feasible region
usually comes from the decision-maker's environment putting some restrictions/conditions on achieving his/her objective.
The following is a very simple illustrative problem. However, the way we approach the problem is the same for a wide variety of decision-making problems, and the
size and complexity may differ. The first example is a product-mix problem.

The Carpenter's Problem:

Allocating Scarce Resources Among Competitive Means
During a couple of brain-storming sessions with a carpenter (our client), he told us that he, solely, makes tables and chairs, sells all tables and chairs he makes at a
market place, however, does not have a stable income, and wishes to do his best.
The objective is to find out how many tables and chairs he should make to maximize net income. We begin by focusing on a time frame, i.e.,
to revise our solution weekly if needed. To learn more about his problem, we must go to his shop and observe what is going on and measure
formulate (i.e., to give a Form, to make a model) of his problem. We must confirm that his objective is to maximize net income. We must communicate with the client.
The carpenter's problem deals with finding out how many tables and chairs to make per week; but first an objective function must be established:
Since the total cost is the sum of the fixed cost (F) and the variable cost per unit multiplied by the number of units produced. Therefore, the decision problem is to find
X1 and X2 such that:
Maximize 9X1 + 6X2 – [(1.5X1 + X2) + (2.5X1 + 2X2) + F1 + F2],
where X1 and X2 stand for the number of tables and chairs; the cost terms in the brackets are the raw material, and labor costs respectively. F1 and F2 are the fixed
costs for the two products respectively. Without loss of generality, and any impact on optimal solution, let us set F1 = 0, and F2 = 0. The objective function reduces to
the following net profit function:
Maximize 5X1 + 3X2
That is, the net incomes (say, in dollars, or tens of dollars) from selling X1 tables and X2 chairs.
The constraining factors which, usually come from outside, are the limitations on labors (this limitation comes from his family) and raw material resources (this
limitation comes from scheduled delivery). Production times required for a table and a chair are measured at different times of day, and estimated to be 2
hour, respectively. Total labor hours per week are only 40 hrs. Raw materials required for a table and a chair are 1, and 2 units respectively. Total supply of raw
material is 50 units per week. Therefore, the LP formulation is:
Maximize 5 X1 + 3 X2
Subject to:
2 X1 + X2  40 labor constraint
X1 + 2 X2  50 material constraint
and both X1, X2 are non-negative.
This is a mathematical model for the carpenter's problem. The decision variables, i.e., controllable inputs are X1, and X2. The output for this model is the total net
income 5 X1 + 3 X2. All functions used in this model are linear (the decision variable have power equal to 1). The coefficients of these constraints are called
Technological Factors (matrix). The review period is one week, an appropriate period within which the uncontrollable inputs (all parameters such as 5, 50, 2,..) are
less likely to change (fluctuate). Even for such a short planning time-horizon, we must perform the what-if analysis to react to any changes in these inputs in order to
control the problem, i.e., to update the prescribed solution.
Notice that since the carpenter is not going out of business at the end of the planning horizon, we added the conditions that both X1, X2 must be non-negative
instead of the requirements that X1, and X2 must be positive integers. The non-negativity conditions are also known as "implied constraints." Again, a Linear
Program would be fine for this problem if the carpenter were going to continue to manufacture these products. The partial items would simply be counted as work in
progress and would eventually become finished goods say, in the next week.
We may try to solve for X1 and X2 by listing possible solutions for each and selecting the pair (X1, X2) that maximize 5X1 + 3X2 (the net income). However, it is too
time consuming to list all possible alternatives and if the alternatives are not exhaustively listed, we cannot be sure that the pair we select (as a solution) is the best
of all alternatives. This way of solving a problem is known as "sequential thinking" versus "simultaneous thinking". More efficient and effective methodologies, known
as the Linear Programming Solution Techniques are based on simultaneous thinking are commercially available in over 400 different software packages from all over
the world.
The optimal solution, i.e., optimal strategy, is to make X1 = 10 tables, and X2 = 20 chairs. We may program the carpenter's weekly activities to make 10 tables and
20 chairs. With this (optimal) strategy, the net income is $110. This prescribed solution was a surprise for the carpenter since, because of more net income of selling
a table ($5), he used to make more tables than chairs!
Hire or Not? Suppose the carpenter can hire someone to help at a cost of $2 per hour. This is, in addition, hourly-based wage he/she is currently paying; otherwise
$2 is much lower than the current minimum wage in US. Should the carpenter hire and if yes then for how may hours?
Let X3 be the number of extra hours, then the modified problem is:
Maximize 5 X1 + 3 X2 - 2 X3
Subject to:
2 X1 + X2  40 + X3 labor constraint with unknown additional hours
X1 + 2 X2  50 material constraint
Under this new condition, we will see that the optimal solution is X1 = 50, X2 = 0, X3 = 60, with optimal net income of $130. Therefore, the carpenter should be hired
for 60 hours. What about only hiring 40 hours? The answer to this and other types of what-if questions are treated under sensitivity analysis in this Web site.

As an exercise, use your LP software to find the largest range for X values satisfying the following inequality with two absolute value terms:
| 3X – 4 | - | 2X – 1 |  2

12 Equation in 3-D space

Parametric vector equations

You can find the vector equation of any straight line through 3D space, but suppose a mathematical model you were working on required the vector equation of a
general curve through 3D space - such as a roller coaster, or trajectory of a satellite.
A parametric curve is defined by a function. The function can be given any real number, known as the parameter t (from a chosen segment of the real number line [a,
b]) and returns a vector, the point on the line corresponding to t. That is, it is a mapping from the segment [a,b] -> R3 3D space.
A convenient way of writing this mapping is with three normal functions, each a mapping from [a,b] -> R:
x = f(t)
y = g(t)
z = h(t)
for each of the cartesian components of the vector returned by the earlier mapping. Here are some examples...

Parametric equation of a circle

From the diagram it can be seen that, if we choose t from the interval [0, 2p] and map t to the vector
(cos t, sin t, 0) we get the unit circle. The separate functions are:
x = cos t
y = sin t
z=0
What curve is given by adding coefficients,
x = a cos t
y = b sin t
z=0
for some real numbers a and b?
Parametric equation of a cylindrical spiral

The parameter, t, can be thought of as time, and the unit circle above is then traced out by a point which starts at (1,0,0) at t = 0 and follows the circular path
counterclockwise (looking down the z axis towards -ve inf.) until it arrives back there at t = 2p. Now if the z function is changed from z = 0 to z = t, the point will rise
steadily as the circular path is traversed, so creating the sprial shown.

Thinking of the curve as being traced out by a particle with the parameter t representing the time is used in Newtonian mechanics. The parametric functions above
give the position of the particle at any time t, and if you have started calculus, differentiating each of the three equations (for x, y and z) with respect t, results in a
new parametric curve - the positions of the points on this curve giving the velocity of the particle at any time t. Differentiating again give the acceleration at any time t.

Parametric equation of a cardiod

x = (cos t)(1 - sin t)

y = (sin t)(1 - sin t)
z=0

Parametric equation of Archimedes' spiral

The spiral found by Archimedes, and also found in nature:
x = t cos t
y = t sin t
z=0

Parametric equation of the Lissajous curves

The curves of Lissajous are the set with equations:

x = cos pt
y = sin qt
z=0
for any p, q belonging to the integers. (the curve on the left has p = 3, q = 5)

Some more complicated parametric equations

The butterfly fly curve, discovered by Temple H. Fay, is generated by the equations
x = (cos t)(ecos t - 2cos 4t - sin5(t/12))
y = (sin t)(ecos t - 2cos 4t - sin5(t/12))
z=0
for t belonging to [0, 24p]

Just as a note, this form of equations isn't anything different to what was covered earlier. For example the vector equation of a line was written,
r(t) = a + t (b - a )
which is the same as,
x = a1(1 - t) + b1t
y = a2(1 - t) + b2t
z = a3(1 - t) + b3t

13. Quadratic equation

Finding square roots by using a calculator is easy for us, but was more of a problem for the Babylonians. In fact they developed a method of successive
approximation to the answer which is identical to the algorithm (called the Newton-Raphson method) used by modern computers to solve much harder problems than
quadratic equations.

Now, not all fields are square. Let's now suppose that the farmer has a more oddly shaped field with two triangular sections as shown on the right.
For appropriate values of a and b the amount of crop that the farmer can grow in this field is given by

This looks a lot more like the quadratic equation that we are used to, and even under the evil eye of the tax man, it’s a lot harder to solve. Yet the Babylonians came
up with the answer again. First we divide by a to give

Now we complete the square by using the fact that

Combining this with the original equation we have

This is now an equation that we can solve by taking square roots. The result is the famous “-b formula":

which can be rewritten as

(The formula usually has "-4ac" because the quadratic equation is more usually written in the form "
The fact that taking a square root can give a positive or a negative answer leads to the remarkable result that a quadratic equation has
mathematical puzzles only having one solution!
Now, this is where the teaching of quadratic equations often stops. We have reached that object beloved of all journalists when they interview mathematicians - a
formula. Endless questions can be made up which involve putting values of a and b into the formula to give (two) answers. But this isn’t what mathematics is about at
all. Finding a formula is only the first step on a long road. We have to ask, what does the formula mean; what does it tell us about the universe; does having a
formula really matter? Let’s now see where this formula will take us.

A quadratic equation is a trinomial of the form ax2 + bx + c = 0. There are three main ways of solving quadratic equations, that are covered below.
Factoring
If a quadratic equation can be factored, then it can be written as a product of two binomials. Since the trinomial is equal to 0, one of the two binomial factors must also
be equal to zero. By factoring the quadratic equation, we can equate each binomial expression to zero and solve each for x.
Examples

Completing the Square

Quadratic equations cannot always be solved by factoring. They can always be solved by the method of completing the squares. To complete the square means to
convert a quadratic to its standard form. We want to convert ax2+bx+c = 0 to a statement of the form a(x - h)2 + k = 0. To do this, we would perform the following
steps:
1) Group together the ax2 and bx terms in parentheses and factor out the coefficient a.

2) In the parentheses, add and subtract (b/2a)2, which is half of the x coefficient, squared.

3) Remove the term - (b/2a)2 from the parantheses. Don't forget to multiply the term by a, when removing from parentheses.

4) Factor the trinomial in parentheses to its perfect square form, (x + b/2a)2.

Note: If you group the - (b/2a)2 + c terms together in parentheses, the equation will now be in standard form. The equation is now much simpler to graph as you will
see in the Graphing section below. To solve the quadratic equation, continue the following steps.
5) Transpose (or shift) all other terms to the other side of the equation and divide each side by the constant a.
6) Take the square root of each side of the equation.

7) Transpose the term -b/2a to the other side of the equation, isolating x.

The quadratic equation is now solved for x. The method of completing the square seems complicated since we are using variables a,b and c. The examples below
show use numerical coefficients and show how easy it can be.
Examples

Note: For more examples of solving a quadratic equation by completing the square, see questions #1 and #2 in the Additional Examples section at the bottom of the
page.

The Quadratic Formula

The method of completing the square can often involve some very complicated calculations involving fractions. To make calculations simpler, a general formula for
solving quadratic equations, known as the quadratic formula, was derived. To solve quadratic equations of the form ax2 + bx + c = 0, substitute the coefficients a,b
and c into the quadratic formula.

Derivation of the Quadratic Formula

The value contained in the square root of the quadratic formula is called the discriminant. If,

Examples
Graphing Quadratic Equations

The graphical representation of quadratic equations are based on the graph of a parabola. A
parabola is an equation of the form y = a x2 + bx + c. The most general parabola, shown at the right,
has the equation y = x2.
The coefficent, a, before the x2 term determines the direction and the size of the parabola. For values
of a > 0, the parabola opens upward while for values of a < 0, the parabola opens downwards. The
graph at the right also shows the relationship between the value of a and the graph of the parabola.

The vertex is the maximum point for

parabolas with a < 0 or minimum
point for parabolas with a > 0. For
parabolas of the form y = ax2, the
vertex is (0,0). The vertex of a
parabola can be shifted however,
and this change is reflected in the
standard equation for parabolas.
Given a parabola y=ax2+bx+c, we
can find the x-coordinate of the
vertex of the parabola using the
formula x=-b/2a. The standard
equation has the form y = a(x - h)2 + k. The parabola y = ax2 is shifted h units to the right and k units
upwards, resulting in a parabola with vertex (h,k).
Note: You may recognize this equation from the completing the square section above. The method of
completing the square should be used to convert a parabola of general form to its standard form.